Diffusion
Collection
Stable Diffusion Models
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6 items
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Updated
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Pony Diffusion XL is a latent text-to-image diffusion model capable of generating images of Horses, mainly, and other things. For more information about how Stable Diffusion functions, please have a look at 🤗's Stable Diffusion blog.
You can use this with the 🧨Diffusers library from Hugging Face.
from diffusers import StableDiffusionXLPipeline
import torch
pipeline = StableDiffusionXLPipeline.from_pretrained("nroggendorff/ponyxl").to("cuda")
image = pipeline(prompt="a chibi doll").images[0]
image.save("horse.png")
train_batch_size
: 1gradient_accumulation_steps
: 4learning_rate
: 1e-2lr_warmup_steps
: 500mixed_precision
: "fp16"eval_metric
: "mean_squared_error"This model card was written by Noa Roggendorff and is based on the Stable Diffusion v1-5 Model Card.