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Physics Study Guide/Periodic Motion. Springs. Hooke's law. formula_1 is the spring constant and formula_2 is the displacement (the force is the restoring force) For this equation to be approximately accurate, formula_2 needs to be below the spring's elastic limit. If formula_2 is more than the spring's elastic limit, the spring will exhibit "plastic behaviour" meaning that it will not return to its original position. Potential energy. potential energy stored in a compressed recoiled spring formula_5 is the period, formula_6 is the mass, and formula_1 is the spring constant.

Norwegian/Lesson 3. Lesson 3 ~ "Kjekt å vite" In this lesson the plan is to elaborate on what you learnt in Lesson 1, in other words, to get a more useful conversation going plus teaching you some things that are simply good to know. Where Are You From? Here are some guys at a course to learn Norwegian. Each person says who they are, what nationality they are, and what languages they speak. Aleksander is the teacher. Check out the vocabulary to understand what everybody is saying. Sorry, I Don't Understand. David is an Englishman. He knows some Norwegian, but has forgotten what "Hvordan går det" means. Are You French? Here are four short exchanges about nationality. May I Introduce Myself? The most common way Norwegians introduce themselves are by simply saying hello, and then telling you their name. For more formal situations, it's more common to greet by first saying your full name, and possibly a more proper greeting than just hello. Hyggelig å treffe deg or Hyggelig å treffe Dem means "pleased to meet you". The first time you meet somebody, you shake hands. In business environments, this is also done when a long amount of time has passed since you last met. When you want to introduce two people to each other, you use the phrase Får jeg presentere ... which means May I present ... and then you say their names. This introduction is quite formal. In less formal situations, you can basically just say Dette er means "this is". It's also not really necessary to use surnames in an informal situation. Especially so with young people and children. Greetings. God morgen, "good morning", is the first greeting of the day. You answer it simply with god morgen or morn', the latter being sort of a slang. Later on in the day you can say god dag, "good day", answered by a god dag back. Being less formal, you can just say hei or hallo; these work at any time of the day. Later, when the evening has come, you say god kveld, which means "good evening", or you could say god aften, which is a synonymous (yet more formal) to god kveld. Finally, when you're going to bed, it's common to say god natt, "good night" or natta', which is a slang expression for the same. When you're leaving somebody, you can say ha det bra or just ha det (lit. "have it", which is short for ha det bra, which implies jeg håper du vil ha det bra, which basically means "I hope you'll be fine". This is a very old saying, and as you can guess, "have it" wouldn't make much sense to anyone if it hadn't been for the history of the word, which everybody sort of knows, but nobody speaks of.) You can also just say snakkes, (short for vi snakkes senere, "we'll talk later") although this is more commonly used on the phone or in communications over written media. Thank You. Takk means "thanks". Tusen takk (lit. "thousand thanks"), Takker så mye/meget (lit. "thank you so much") and Mange takk (lit. "many thanks") all means "thanks a lot". Using meget is old-fashioned though, and comes from Danish. The correct response for takk, for instance when you've given someone a gift, or similar is vær så god (lit. "be so kind") which is equivalent to the English "you are welcome", "don't mention it" or "it's nothing" When someone says takk in order to express gratitude for a service or act performed, bare hyggelig, (lit. "merely pleasurable"), best translated as "my pleasure", is an appropriate response.

XML - Managing Data Exchange/Data Schemas (Answers). XML = Exercises = Answers by Danny Popov Question 1 chapter4.xsd Chap4Redefined.xsd Chap4Redefined.xml Question 2 USAddress.xsd Bar.xsd Bar.xml Question 3 Museum.xsd Museum.xml Question 4 TypeLibrary.xsd

Norwegian. Vocabulary. /Months and Weekdays/ | /Numbers/ | /Time/ | /Useful Expressions/ Appendices. Advanced Noun Inflection | Norwegian Pronunciation | Misunderstandings | Fast translation rules See also. __NOEDITSECTION__

XML - Managing Data Exchange/XSLT and Style Sheets. In previous chapters, we have introduced the basics of using an XSL stylesheet to convert XML documents into HTML. This chapter will briefly review those concepts and introduce many new ones as well. It is a reference for creating stylesheets. XML Stylesheets. The eXtensible Stylesheet Language (XSL) provides a means to transform and format the contents of XML document for display. It includes two parts, XSL Transformation (XSLT) for transforming the XML document, and XSLFO (XSL Formatting Objects) for formatting or applying styles to XML documents. The XSL Transformation Language (XSLT) is used to transform XML documents from one form to another, including new XML documents, HTML, XHTML, and text documents. XSL-FO can create PDF documents, as well as other output formats, from XML. With XSLT you can effectively recycle content, redesigning it for use in new documents, or changing it to fit limitless uses. For example, from a single XML source file, you could extract a document ready for print, one for the Web, one for a Unix manual page, and another for an online help system. You can also choose to extract only parts of a document written in a specific language from an XML source that stores text in many languages. The possibilities are endless! An XSLT stylesheet is an XML document, complete with elements and attributes. It has two kinds of elements, top-level and instruction. Top-level elements fall directly under the codice_1 root element. Instruction elements represent a set of formatting instructions that dictate how the contents of an XML document will be transformed. During the transformation process, XSLT analyzes the XML document, or the source tree, and converts it into a node tree, a hierarchical representation of the entire XML document, also known as the result tree. Each node represents a piece of the XML document, such as an element, attribute or some text content. The XSL stylesheet contains predefined “templates” that contain instructions on what to do with the nodes. XSLT will use the codice_2 attribute to relate XML element nodes to the templates, and transform them into the result document. Let's review the stylesheet, city.xsl from chapter 2, and examine it in a little more detail: Exhibit 1: XML stylesheet for city entity <?xml version="1.0" encoding="UTF-8"?> Document: city.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="html"/> <xsl:template match="/"> <html> <head> <title>Cities</title> </head> <body> <h2>Cities</h2> <xsl:apply-templates select="cities"/> </body> </html> </xsl:template> <xsl:template match="cities"> <!-- the for-each element can be used to loop through each node in a specified node set (in this case city) --> <xsl:for-each select="city"> <xsl:text>City: </xsl:text> <xsl:value-of select="cityName"/> <br/> <xsl:text>Population: </xsl:text> <xsl:value-of select="cityPop"/> <br/> <xsl:text>Country: </xsl:text> <xsl:value-of select="cityCountry"/> <br/> <br/> </xsl:for-each> </xsl:template> </xsl:stylesheet> The codice_3 element defines the rules that implement a change. This can be any number of things, including a simple plain-text conversion, the addition or removal of XML elements, or simply a conversion to HTML, when the pattern is matched. The pattern, defined in the element’s codice_2 attribute, contains an abbreviated ../XPath/ location path. This is basically the name of the root element in the doc, in our case, "tourGuide." When transforming an XML document into HTML, the processor expects that elements in the stylesheet be well-formed, just as with XML. This means that all elements must have an end tag. For example, it is not unusual to see the codice_5 tag alone. The XSLT processor requires that an element with a start-tag must close with an end tag. With the codice_6 element, this means either using codice_7 or codice_8 As mentioned in Chapter 3, the codice_9 element is an empty element. That means it carries no content between tags, but it may have attributes. Although no end tags are output for the HTML output, they still must have end-tags in the stylesheet. For instance, in the stylesheet, you will list: codice_10 or as an empty element codice_11. The HTML output will drop the end-tag so it looks like this: codice_12 On a side note, the processor will recognize html tags no matter what case they are in - BODY, body, Body are all interpreted the same. Output. XSLT can be used to transform an XML source into many different types of documents. XHTML is also XML, if it is well formed, so it could also be used as the source or the result. However, transforming plain HTML into XML won't work unless it is first turned into XHTML so that it conforms to the XML 1.0 recommendation. Here is a list of all the possible type-to-type transformations performed by XSLT: Exhibit 2: Type-To-Type Transformations The codice_13 element in the stylesheet determines how to output the result tree. This element is optional, but it allows you to have more control over the output. If you do not include it, the output method will default to XML, or HTML if the first element in the result tree is the codice_14element. Exhibit 3 lists attributes. Exhibit 3: Element output attributes XML to XML. Since we have had a lot of practice transforming an XML document to HTML, we are going to transform city.xml, used in chapter 2, into another XML file, using host.xsd as the schema. Exhibit 4: XML document for city entity <?xml version="1.0" encoding="UTF-8"?> Document: city.xml <cities xmlns:xsi='http://www.w3.org/2001/XMLSchema-instance' xsi:noNamespaceSchemaLocation='host.xsd'> <city> <cityID>c1</cityID> <cityName>Atlanta</cityName> <cityCountry>USA</cityCountry> <cityPop>4000000</cityPop> <cityHostYr>1996</cityHostYr> </city> <city> <cityID>c2</cityID> <cityName>Sydney</cityName> <cityCountry>Australia</cityCountry> <cityPop>4000000</cityPop> <cityHostYr>2000</cityHostYr> <cityPreviousHost>c1</cityPreviousHost > </city> <city> <cityID>c3</cityID> <cityName>Athens</cityName> <cityCountry>Greece</cityCountry> <cityPop>3500000</cityPop> <cityHostYr>2004</cityHostYr> <cityPreviousHost>c2</cityPreviousHost > </city> </cities> Exhibit 5: XSL document for city entity that list cities by City ID <?xml version="1.0" encoding="UTF-8"?> <xsl:stylesheet xmlns:xsl="http://www.w3.org/1999/XSL/Transform" version="2.0"> <xsl:output method="html"/> <xsl:template match="/"> <xsl:for-each select="//city[count(cityPreviousHost) = 0]"> <br/><xsl:text>City Name: </xsl:text><xsl:value-of select="cityName"/><br/> <xsl:text> Rank: </xsl:text><xsl:value-of select="cityID"/><br/> <xsl:call-template name="output"> <xsl:with-param name="context" select="."/> </xsl:call-template> </xsl:for-each> </xsl:template> <xsl:template name="output"> <xsl:param name="context" select="."/> <xsl:for-each select="//city[cityPreviousHost = $context/cityID]"> <br/><xsl:text>City Name: </xsl:text> <xsl:value-of select="cityName"/><br/> <xsl:text> Rank: </xsl:text><xsl:value-of select="cityID"/><br/> <xsl:call-template name="output"> <xsl:with-param name="context" select="."/> </xsl:call-template> </xsl:for-each> </xsl:template> </xsl:stylesheet> Exhibit 6: XML schema for host city entity <?xml version="1.0" encoding="UTF-8"?> <xsd:schema xmlns:xsd="http://www.w3.org/2001/XMLSchema" elementFormDefault="qualified" attributeFormDefault="unqualified"> <xsd:element name="cities"> <xsd:complexType> <xsd:sequence> <xsd:element name="city" type="cityType" maxOccurs="unbounded"/> </xsd:sequence> </xsd:complexType> </xsd:element> <xsd:complexType name="cityType"> <xsd:sequence> <xsd:element name="cityID" type="xsd:ID"/> <xsd:element name="cityName" type="xsd:string"/> <xsd:element name="cityCountry" type="xsd:string"/> <xsd:element name="cityPop" type="xsd:integer"/> <xsd:element name="cityHostYr" type="xsd:integer"/> <xsd:element name="cityPreviousHost" type="xsd:IDREF" minOccurs="0" maxOccurs="1"/> </xsd:sequence> </xsd:complexType> Exhibit 7: XML stylesheet for city entity <?xml version="1.0" encoding="UTF-8"?> Document: city2.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="xml" encoding="utf-8" indent="yes" /> <xsl:attribute-set name="date"> <xsl:attribute name="year">2004</xsl:attribute> <xsl:attribute name="month">03</xsl:attribute> <xsl:attribute name="day">19</xsl:attribute> </xsl:attribute-set> <xsl:template match="tourGuide"> <xsl:processing-instruction name="xsl-stylesheet"> href="style.css" type="text/css"<br /> </xsl:processing-instruction> <xsl:comment>This is a list of the cities we are visiting this week</xsl:comment> <xsl:for-each select="city"> <!-- element name creates a new element where the value of the attribute name sets name of the new element. Multiple attribute sets can be used in the same element --> <!-- use-attribute-sets attribute adds all the attributes declared in attribute-set from above --> <xsl:element name="cityList" use-attribute-sets="date"> <xsl:element name="city"> <xsl:attribute name="country"> <xsl:apply-templates select="country"/> </xsl:attribute> <xsl:apply-templates select="cityName"/> </xsl:element> <xsl:element name="details">Will write up a one page report of the trip</xsl:element> </xsl:element> </xsl:for-each> </xsl:template> </xsl:stylesheet> The stylesheet produces this result tree: Exhibit 8: XML result tree for city entity <?xml version="1.0" encoding="utf-8" standalone="no"?> Document: city2.xsl <?xsl-stylesheet href="style.css" type="text/css"?> <!--This is a list of the cities we are visiting this week--> <cityList year="2004" month="03" day="19"> <city country="Belize">Belmopan</city> <details>Will write up a one page report of the trip</details> </cityList> <cityList year="2004" month="03" day="19"> <city country="Malaysia">Kuala Lumpur</city> <details>Will write up a one page report of the trip</details> </cityList> </stylesheet> The processor automatically inserts the XML declaration at the top of the result tree. The processing instruction, or PI, is an instruction intended for use by a processing application. In this case, the href points to a local stylesheet that will be applied to the XML document when it is processed. We used codice_22 to create new content in the result tree and added attributes to it. There are two other instruction elements for inserting nodes into a result tree. These are codice_23 and codice_24. Unlike codice_25, which only copies content of the child node (like the child text node), these elements copy everything. The following code shows how the copy element can be used to copy the city element in city.xml: Exhibit 9: Copy element <xsl:template match="city"> <xsl:copy /> </xsl:template> The result looks like this: Exhibit 10: Copy element result <?xml version="1.0" encoding="utf-8"> <city /> <city /> The output isn't very interesting, because copy does not pick up the child nodes, only the current node. In our example, it picks up the two city nodes that are in the city.xml file. The copy element has an optional attribute, use-attribute-sets, which allows you to add attributes to the element. However, it will leave behind any other attributes, except the namespace, if it is present. Here is the result if a namespace is declared in the source document, in this case, the default namespace: Exhibit 11: Namespace result <?xml version="1.0" encoding="utf-8"> <city xmlns:xsi="http://www.w3.org/2001/XMLSchema-instance"> <city xmlns:xsi="http://www.w3.org/2001/XMLSchema-instance"> If you want to copy more from the source file than just one node, the codice_24 element includes the current node, and any attribute nodes that are associated with it. This includes any nodes that might be laying around, such as namespace nodes, text nodes, and child element nodes. When we apply the codice_24 element to city.xml, the result is almost an exact replica of city.xml! You can also copy comments and processing instructions using codice_28 and codice_29where codice_30 is the value of the name attribute in the processing instruction you wish to retrieve. Why would this be useful, you ask? Sometimes you want to just grab nodes and go! For example, if you want to place a copy of city.xml into a SOAP envelope, you can easily do it using codice_24. If you don't already know, Simple Object Access Protocol, or SOAP, is a protocol for packaging XML documents for exchange. This is really useful in a B2B environment because it provides a standard way to package XML messages. You can read more about SOAP at www.w3.org/tr/soap. Use an XML editor to create the above XML Stylesheets, and experiment with the codice_23 and codice_24 elements. Templates. Since templates define the rules for changing nodes, it would make sense to reuse them, either in the same stylesheet or in other stylesheets. This can be accomplished by naming a template, and then calling it with a codice_34 element. Named templates from other stylesheets can also be included. You can quickly see how this is useful in practical applications. Here is an example using named templates: Exhibit 110: Named templates <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="xml" /> <xsl:template match=" /"> <xsl:call-template name="getCity" /> </xsl:template> <xsl:template name="getCity"> <xsl:copy-of select="city" /> </xsl:template> </xsl:stylesheet> Templates also have a mode attribute. This allows you to process a node more than once, producing a different result each time, depending on the template. Let's create a stylesheet to practice modes. Exhibit 12: XML template modes <?xml version="1.0" encoding="UTF-8"?> Document: cityModes.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="html" /> <xsl:template match="tourGuide"> <html> <head> <title>City - Using Modes</title> </head> <body> <xsl:for-each select="city"> <xsl:apply-templates select="cityName" mode="title" /> <xsl:apply-templates select="cityName" mode="url" /> <br /> </xsl:for-each> </body> </html> </xsl:template> <xsl:template match="cityName" mode="title"> <h2><xsl:value-of select="current()"/></h2> </xsl:template> <xsl:template match="cityName" mode="message"> <p>Come visit <b><xsl:value-of select="current()" /></b>!</p> </xsl:template> </xsl:stylesheet> The result isn't very flattering since we didn't do much with the file, but it gets the point across. Exhibit 13: Result from above stylesheet <h2>Belmopan</h2> Come visit <b>Belmopan</b>! <h2>Kuala Lumpur</h2> Come visit <b>Kuala Lumpur</b>! By default, XSLT processors have built-in template rules. If you apply a stylesheet without any matching rules, and it fails to match a pattern, the default rules are automatically applied. The default rules output the content of all the elements. Sorting. Writing “well formed” code XML is vital. At times, however, simply displaying information (the most elementary level of data management) is not all that is necessary to properly identify a project. As information technology specialists, it is necessary to fully understand that order is vital for interpretation. Order can be attained by putting data in a format that is quickly readable. Such information then becomes quickly usable. Using a comparative model or simply looking for a specific name or item becomes very easy. Finding a specific musical artist, title, or musical type becomes very easy. As an Information Specialist, you must fully be aware that it often becomes necessary to sort information. The basis of sorting in XMLT is the xsl:sort command. The xsl:sort element exemplifies a sort key component. A sort key component identifies how a sort key value is to be identified for each item in the order of information being sorted. A Sort Key Value is defined as “the value computed for an item by using the Nth sort key component” The significance of a sort key component is realized either by its select attribute, or by the contained sequence constructor. A Sequence Constructor is defined as a “sequence of zero or more sibling nodes in the stylesheet that can be evaluated to return a sequence of nodes and atomic values”. There are instances when neither is present. Under these circumstances, the default is select=".", which has the effect of sorting on the actual value of the item if it is an atomic value, or on the typed-value of the item if it is a node. If a select attribute is present, its value must be an Xpath expression. The following is how the <xsl:sort> element is used to sort the output. Sort Information is held as Follows: Sorting output in XML is quite easy and is done by adding the <xsl:sort> element after the <xsl:for-each> element in the XSL file. Exhibit 14: Stylesheet with sort function <?xml version="1.0" encoding="ISO-8859-1"?> <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:template match="/"> <html> <body> <h2>TourGuide Example</h2> <xsl:apply-templates select="cities"/> </body> </html> </xsl:template> <xsl:template match="cities"> <xsl:for-each select="city"> <xsl:sort select="cityName"/> <xsl:value-of select="cityName"/> <xsl:value-of select="cityCountry"/> </xsl:for-each> </xsl:template> </xsl:stylesheet> This example will sort the file alphabetically by artist name. Note: The select attribute indicates what XML element to sort on. Information can be SELECTED and SORTED by “title” or “artist”. These are categories that the XML document will display within the body of the file. We have used the codice_39 function to sort the results of an codice_40 statement before. The sort element has many other uses as well. Essentially, it instructs the processor to sort nodes based on certain criteria, which is known as the sort key. It defaults to sorting the elements in ascending order. Here is a short list of the different attributes that sort takes: Exhibit 15: Sort attributes The sort element can be used in either the codice_25 or the codice_42 elements. It can also be used multiple times within a template, or in several templates, to create sub-ordering levels. Numbering. The codice_43 instruction element allows you to insert numbers into your results. Combined with a sort element, you can easily create numbered lists. When this simple stylesheet, hotelNumbering.xsl, is applied to city_hotel.xml, we get the result listed below: Exhibit 16: Sorting and numbering lists <?xml version="1.0" encoding="ISO-8859-1"?> Document: hotelNumbering.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="text" omit-xml-declaration="yes"/> <xsl:template match="/"> <xsl:apply-templates select="tourGuide/city/hotel"> <xsl:sort/> </xsl:apply-templates> </xsl:template> <xsl:template match="hotel"> <xsl:number value="position()" format="
 0. "/> <xsl:value-of select="hotelName"/> </xsl:template> </xsl:stylesheet> Exhibit 17: Result hotelNumbering.xsl 1. Bull Frog Inn 2. Mandarin Oriental Kuala Lumpur 3. Pan Pacific Kuala Lumpur 4. Pook's Hill Lodge The expression in codice_44 is evaluated and the value for codice_45 is based on the sorted node list. To improve the looks we are adding the format attribute with a linefeed character reference (
), a zero digit to indicate that the number will be a zero digit to indicate that the number will be an integer type, and a period and space to make it look nicer. The format list can be based on the following sequences: Exhibit 17: Numbering formats codice_46 – Uppercase letters codice_47 – Lowercase letters codice_48 – Uppercase Roman numerals codice_49 – Lowercase Roman numerals codice_50 – Numeral prefix codice_51 – Integer prefix/ hyphen prefix To specify different levels of numbering, such as sections and subsections of the source document, the codice_52 attribute is used, which tells the processor the levels of the source tree that should be considered. By default, it is set to codice_53, as seen in the example above. It also can take values of codice_54 and codice_55. The codice_56 attribute is a pattern that tells the processor which nodes to count (for numbering purposes). If it is not specified, it defaults to a pattern matching the same node type as the current node. The codice_57 attribute can also be used to specify the node where the counting should start. When level is set to codice_53, the processor searches for nodes that match the value of codice_56, and if it is not present, it matches the current node. When it finds the match, it creates a node-list and counts all the matching nodes of that type. If the codice_57 attribute is listed, it tells the processor where to start counting from, rather than counting all nodes When the level is codice_54, it doesn't just count a list of one node type, it creates a list of all the nodes that are ancestors of the current node, in the actual order from the source document. After this list is created, it selects all the nodes that match the nodes represented in count. It then maps the number of preceding siblings for each node that matches count. In effect, codice_54 remembers all the nodes separately. This is where codice_55 is different. It will number all the elements sequentially, instead of counting them in multiple levels. As with the other two values, you can use the codice_57 attribute to tell the processor where to start counting from, which in effect will separate it into levels. This is a modification of the example above using the codice_65: Exhibit 18: Sorting and numbering lists Document: hotelNumbering2.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="text" omit-xml-declaration="yes"/> <xsl:template match="/"> <xsl:apply-templates select="tourGuide//hotelName"/> </xsl:template> <xsl:template match="hotel"> <xsl:number level="multiple" count="city|hotel" format="
 1.1 "/> <xsl:apply-templates /> </xsl:template> </xsl:stylesheet> Exhibit 19: Result – hotelNumbering2.xsl 1.1 Bull Frog Inn 1.2 Pook's Hill Lodge 2.1 Pan Pacific Kuala Lumpur 2.2 Mandarin Oriental Kuala Lumpur The first template matches the root node and then selects all codice_66 nodes that have codice_67 as an ancestor, creating a node-list. The next template recursively processes the codice_68 element, and gives it a number for each instance of codice_68 based on the number of elements in the attribute. This is figured out by counting the number of preceding siblings, plus 1. Formatting. Formatting numbers is a simple process so this section will be a brief overview of what can be done. Placed within the XML stylesheet, functions can be used to manipulate data during the transformation. In order to make numbers a little easier to read, we need to be able to separate the digits into groups, or add commas or decimals. To do this we use the codice_70 function. The purpose of this function is to convert a numeric value into a string using specified patterns that control the number of leading zeroes, separator between thousands, etc. The basic syntax of this function is as follows: codice_71 The following are the characters and their meanings used to represent the number format when using the format-number function within a stylesheet: Exhibit 20: Format-number function Symbol Meaning 0 A digit. # A digit, zero shows as absent. . (period) Placeholder for decimal separator. , Placeholder for grouping separator. ; Separate formats. - Default prefix for negative. % Multiply by 100 and show as a percentage. X Any other characters can be used in the prefix or suffix. ‘ Used to quote special characters in a prefix or suffix. Conditional Processing. There are times when it is necessary to display output based on a condition. There are two instruction elements that let you conditionally determine which template will be used based on certain tests. These are the codice_40 and codice_75 elements. The test condition for an codice_40 statement must be contained within the codice_77 attribute of the codice_78 element. Expressions that are testing greater than and less than operators must represent them by “>” and “<” respectively in order for the appropriate transformation to take place. The codice_79 function from XPath is a Boolean function and evaluates to true if its argument is false, and vice versa. The codice_80 and codice_81 conditions can be used to combine multiple tests, but an codice_40 statement can, at most, test only one expression. It can also only instantiate the use of one template. The codice_83 element, is similar to the codice_84 statement in Java. By using the codice_83 element, the codice_75 element can offer a many alternative expressions. A choose element must contain at least one when statement, but it can have as many as it needs. The choose element can also contain one instance of the otherwise element, which works like the final else in a Java program. It contains the template if none of the other expressions are true. The codice_42 element is another conditional processing element. We have used it in previous chapter exercises, so this will be a quick review. The codice_42 element is an instruction element, which means it must be children of template elements. codice_42 evaluates to a node-set, based on the value of the select attribute, or expression, and processes through each node in document order, or sorted order. Parameters and Variables. XSLT offers two similar elements, codice_90 and codice_91. Both have a required codice_30 attribute, and an optional codice_93 attribute, and you declare them like this: Exhibit 21: Variable and parameter declaration <xsl:variable name="var1" select="""/> <xsl:param name="par1" select="""/> The above declarations have bound to an empty string, which is the same effect as if you had left off the select attribute. With parameters, this value is considered only a default, or initial value to be changed either from the command line, or from another template using the codice_94 element. However, with the variable, as a general rule, the value is set and can't be changed dynamically except under special circumstances. When making declarations, remember that variables can be declared anywhere within a template, but a parameter must be declared at the beginning of the template. Both elements can also have global and local scope, depending on where they are defined. If they are defined at the top-level under the <stylesheet> elements, they are global in scope and can be used anywhere in the stylesheet. If they are defined in a template, they are local and can only be used in that template. Variables and parameters declared in templates are visible only to the template they are declared in, and to templates underneath them. They have a cascading effect: they can spill down from the top-level into a template, down into a template within that one, etc, but they cannot go back up! We are going to hard-code a value for the parameter in it's declaration element using the codice_93 attribute. Exhibit 22: HTML results <?xml version="1.0" encoding="UTF-8"?> Document: countryParam.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="text"/> <xsl:param name="country" select="'Belize'"/> <xsl:param name="code" /> <xsl:template match="/"> <xsl:apply-templates select="country-codes" /> </xsl:template> <xsl:template match="country-codes"> <xsl:apply-templates select="code" /> </xsl:template> <xsl:template match="code"> <xsl:choose> <xsl:when test="countryName[. = $country]"> The country code for <xsl:value-of select="countryName"/> is <xsl:value-of select="countryCode"/>. </xsl:when> <xsl:when test="countryCode[. = $code]"> The country for the code <xsl:value-of select="countryCode"/> is <xsl:value-of select="countryName"/>. </xsl:when> <xsl:otherwise> Sorry. No matching country name or country code. </xsl:otherwise> </xsl:choose> </xsl:template> </xsl:stylesheet> The value that you pass in does not have to be enclosed in quotes, unless you are passing a value with more than one word. For example, we could have passed either country="United States" or country=Belize without getting an error. The value of a variable can also be used to set an attribute value. Here is an example setting the countryName element with an attribute of countryCode equal to the value in thecodice_96 variable: Exhibit 23: Attribute of countryCode codice_97 This is known as an attribute value template. Notice the use of braces around the parameter. This tells the processor to evaluate the content as an expression, which then converts the result to a string in the result tree. There are attributes which cannot be set with an attribute value template: Parameters, though not variables, can be passed between templates using the codice_94 element. This element has two attributes, codice_30, which is required, and codice_93, which is optional. This next example uses with-param as a child of the codice_34 element, although it can also be used as a child of codice_25. Exhibit 24: XSL With-Param <?xml version="1.0" encoding="UTF-8"?> Document: withParam.xsl <xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"> <xsl:output method="text"/> <xsl:template match="/"> <xsl:apply-templates select="tourGuide/city"/> </xsl:template> <xsl:template match="city"> <xsl:call-template name="countHotels"> <xsl:with-param name="num" select="count(hotel)"/> </xsl:call-template> </xsl:template> <xsl:template name="countHotels"> <xsl:param name="num" select=""" /> <xsl:text>City Name: </xsl:text> <xsl:value-of select="cityName" /> <xsl:text>
</xsl:text> <xsl:text>Number of hotels: </xsl:text> <xsl:value-of select="$num" /> <xsl:text>

</xsl:text> </xsl:template> </xsl:stylesheet> Exhibit 25: Text results – withParam.xsl City Name: Belmopan Number of hotels: 2 City Name: Kuala Lumpur Number of hotels: 2 The Muenchian Method. The Muenchian Method is a method developed by Steve Muench for performing functions using keys. Keys work by assigning a key value to a node and giving you access to that node through the key value. If there are lots of nodes that have the same key value, then all those nodes are retrieved when you use that key value. Effectively this means that if you want to group a set of nodes according to a particular property of the node, then you can use keys to group them together. One of the more common uses for the Muenchian method is grouping items and counting the number of occurrences in that group, such as number of occurrences of a city <?xml version="1.0" encoding="UTF-8"?> <xsl:stylesheet xmlns:xsl="http://www.w3.org/1999/XSL/Transform" version="1.0"> <xsl:output method="html"/> <xsl:key name="Count" match="*/city" use="cityName" /> <xsl:template match="cities"> <xsl:for-each select="//city[generate-id()=generate-id(key('Count', cityName)[1])]"> <br/><xsl:text>City Name:</xsl:text><xsl:value-of select="cityName"/><br/> <xsl:text>Number of Occurences:</xsl:text> <xsl:value-of select="count(key('Count', cityName))"/> <br/> </xsl:for-each> </xsl:template> </xsl:stylesheet> Text Results – muenchianMethod.xsl City Name: Atlanta Number of Occurrences: 1 City Name: Athens Number of Occurrences: 1 City Name: Sydney Number of Occurrences: 1 Datatypes. There are five different datatypes in XSLT: Node-set, String, Number, Boolean, and Result tree fragment. Variables and parameters can be bound to each of these, but the last type is specific to them. Node-sets are returned everywhere in XSLT. We've seen them returned from codice_25 and codice_119 elements, and variables. Now we will see how a variable can be bound to a node-set. Examine the following code: Exhibit 26: Variable bound to a node-set <xsl:variable name="cityNode" select="city" /> <xsl:template match="/"> <xsl:apply-templates select="$cityNode/cityName" /> </xsl:template> Here, we are setting the value of the variable codice_120 to the node-set codice_108 from the source tree. The codice_122 element is a child of city, so the output generated by apply-templates is the text node of codice_122. Remember, you can use variable references in expressions but not patterns. This means we cannot use the reference codice_124 as the value of a codice_2 attribute. String types are useful if you are interested only in the text of nodes, rather than in the whole node-set. String types use XPath functions, most notably, codice_126. This is just a simple example: Exhibit 27: String types <xsl:variable name="cityName" select="string('Belmopan')" /> This is in fact, a longer way of saying: Exhibit 28: Shorter version of above <xsl:variable name="cityName" select="' Belmopan'" /> It is also possible to declare a variable that has a number value. You do this by using the XPath function codice_127. Exhibit 29: Declaration of variable with number value <xsl:variable name="population" select="number(11100)" /> You can use numeric operators such as + - * / to perform mathematic operations on numbers, as well as some built in XPath functions such as codice_128 and codice_129. The Boolean type has only two possible values, true or false. As an example, we are going to use a Boolean variable to test to see if a parameter has been passed into the stylesheet. Exhibit 30: Boolean variable to test <xsl:param name="isOk" select=""" /> <xsl:template match="city" /> <xsl:choose> <xsl:when test="boolean($isOk)"> …logic here… </xsl:when> <xsl:otherwise> Error: must use parameter isOk with any value to apply template </xsl:otherwise> </xsl:choose> </xsl:template> We start with an empty-string declaration for the parameter codice_130. In the codice_77 attribute of codice_83, the codice_133 function tests the value of codice_130. If the value is an empty string, as we defined by default, codice_133 evaluates to codice_136, and the template is not instantiated. If it does have a value, and it can be any value at all, codice_133 evaluates to codice_138. The final datatype is the result tree fragment. Essentially it is a chunk of text (a string) that can contain markup. Let's look at an example before we dive into the details: Exhibit 31: Result tree fragment datatype <xsl:variable name="fragment"> <description>Belmopan is the capital of Belize</description> </xsl:variable> Notice we didn't use the select attribute to define the variable. We aren't selecting a node and getting its value, rather we are creating arbitrary text. Instead, we declared it as the content of the element. The text in between the opening and closing variable tags is the actual fragment of the result tree. In general, if you use the select attribute as we did earlier, and don't specify content when declaring variables, the elements are empty elements. If you don't use select and you do specify content, the content is a result tree. You can perform operations on it as if it were a string, but unlike a node set, you can't use operators such as / or // to get to the nodes. The way you retrieve the content from the variable and get it into the result tree is by using the copy-of element. Let's see how we would do this: Exhibit 32: Retrieve and place into result tree <xsl:template match="city" <xsl:copy-of select="cityName" /> <xsl:copy-of select="$fragment" /> </xsl:template> The result tree would now contain two elements: a copy of the city element and the added element, description. EXSLT. EXSLT is a set of community developed extensions to XSLT. The modules include facilities to handle dates and times, math, and strings. Multiple Stylesheets. In previous chapters, we have imported and used multiple XML and schema documents. It is also possible to use multiple stylesheets using the codice_139 and codice_140 elements, which should be familiar. It is also possible to process multiple XML documents at a time, in one stylesheet, by using the XSLT function codice_141. Including an external stylesheet is very similar to what we have done in earlier chapters with schemas. The codice_140 element only has one attribute, which is codice_143. It is required and always contains a URI (Uniform Resource Identifier) reference to the location of the file, which can be local (in the same local directory) or remote. You can include as many stylesheets as you need, as long as they are at the top level. They can be scattered all over the stylesheet if you want, as long as they are children of the codice_144 element. When the processor encounters an instance of include, it replaces the instance with all the elements from the included document, including template rules and top-level elements, but not the root codice_144 element. All the items just become part of the stylesheet tree itself, and the processor treats them all the same. Here are declarations for including a local and remote stylesheet: Exhibit 33: Declarations for local and remote stylesheet <xsl:include href="city.xsl" /> <xsl:include href="http://www.somelocation.com/city.xsl"/> Since codice_140 returns all the elements in the included stylesheet, you need to make sure that the stylesheet you are including does not include your own stylesheet. For example, city.xsl cannot include city_hotel.xsl, if city_hotel.xsl has an include element which includes city.xsl. When including multiple files, you need to make sure that you are not including another stylesheet multiple times. If city_hotel.xsl includes amenity.xsl, and country.xsl includes amenity.xsl, and city.xsl includes both city_hotel.xsl and country.xsl, it has indirectly included amenity.xsl twice. This could cause template rule duplication and errors. These are some confusing rules, but they are easy to avoid if you carefully examine the stylesheets before they are included. The difference between importing stylesheets and including them is that the template rules imported each have a different import precedence, while included stylesheet templates are merged into one tree and processed normally. Imported templates form an import tree, complete with the root codice_144 element so the processor can track the order in which they were imported. Just like include, import has one attribute, href, which is required and should contain the URI reference for the document. It is also a top-level element and can be used as many times as need. However, it must be the immediate child for the codice_144 element, otherwise there will be errors. This code demonstrates importing a local stylesheet: Exhibit 34: Importing local stylesheet <xsl:import href="city.xsl" /> The order of the codice_139 elements dictates the precedence that matching templates will have over one another. Templates that are imported last have higher priority than those that are imported first. However, the codice_3 element also has a codice_151 attribute that can affect its priority. The higher the number in the codice_151 attribute, the higher the precedence. Import priority only comes into effect when templates collide, otherwise importing stylesheets is not that much different from including them. Another way to handle colliding templates is to use the codice_153 element. If a template in the imported document collides with a template in the importing document, apply-templates will override the rule and cause the imported template to be invoked. The codice_141 function allows you to process additional XML documents and their nodes. The function is called from any attribute that uses an expression, such as the codice_93 attribute. For example: Exhibit 35: Document() function <xsl:template match="hotel"> <xsl:element name="amenityList"> <xsl:copy-of select="document('amenity.xml')" /> </xsl:element> </xsl:template> When applied to an xml document that only contains an empty hotel element, such as codice_156the result tree will add a new element called amenityList, and place all the content from amenity.xml (except the XML declaration) in it. The document function can take many other parameters such as a remote URI, and a node-set, just to name a few. For more information on using codice_141, visit http://www.w3.org/TR/xslt#document XSL-FO. XSL-FO stands for Extensible Stylesheet Language Formatting Objects and is a language for formatting XML data. When it was created, XSL was originally split into two parts, XSL and XSL-FO. Both parts are now formally named XSL. XSL-FO documents define a number of rectangular areas for displaying output. XSL-FO is used for the formatting of XML data for output to screen, paper or other media, such as PDF format. For more information, visit http://www.w3schools.com/xslfo/default.asp Reference Section. Exhibit 36: XSL Elements Exhibit 37: XSLT Functions Exhibit 38: Inherited XPath Functions "Node Set Functions" "String Functions" "Number Functions" " Boolean Functions" Exercises. In order to learn more about XSL and stylesheets, exercises are provided. Answers. In order to learn more about XSL and stylesheets, answers are provided.

Electronics/Electric Field. Electric Field. Charged particles create an electric field. A good way to look at an electric field is electrolysis. cathode = electron emitter. <br> anode = electron collector. (Means up the path). <br> Static Electricity. "This is a rough draft." Generated by contact between dissimilar materials, dissipates through moisture when humid, dry air is a better insulator. "As this is a book, it should have historical stories and stuff to make it more interesting. Greeks and amber in this section, for instance." When a solitary atom has the same number of electrons and protons, it has no overall charge. The individual charges cancel out. It is said to be "neutral". An atom that has more electrons than protons is negatively charged. The atom will have a tendency to attract other particles that have a positive charge, and repel particles with a negative charge. An atom that has less electrons than protons is positively charged. The atom will attract negative particles (like solitary electrons or negatively charged atoms) and repel positive particles (like solitary protons or positively charged atoms). If a solitary proton and electron are placed near each other (or even far away, as long as nothing else is nearby), they will move closer together, until the electron becomes "attached" to the proton and they form a hydrogen atom. An atom of metal that is positively charged (missing electrons) and a free electron will behave the same way. The same goes for macroscopic objects. A metal sphere that has a significant amount of atoms that have the same number of electrons and protons is neutral. Also, if there are several atoms that are positively charged or negatively charged, it is said that the whole object is charged. This may not be noticeable if only a few atoms on the object are charged. Charged metal spheres will attract and repel each other according to the same rules as individual particles. A solitary metal sphere with an excess of electrons will tend to have the extra electrons spread out evenly across the surface of the sphere (not in the interior). This is because the electrons repel each other, and try to get as far apart as possible. On a positively charged sphere, the positive charge (absence of electrons) also spreads evenly across the surface due to its lack of electrons. (Talk about the neat shielding properties of Faraday cages) In the space inside to the sphere there is no charge. In the case of insulating spheres charge does not move throughout the object, it just sits there. The electrons in this object are not constantly moving from atom to atom, so the empty holes have no way to be filled by their neighbors. A Van de Graaff generator is basically a pump for electrons (or a "charge pump"). It uses the triboelectric effect to take electrons from one conductor and deposit them on another. The triboelectric effect is caused when two different materials are brought into contact and then separated. One material will have more of a tendency to attract extra electrons, and some of the electrons will transfer from one material to the other. The Van de Graaff generator accomplishes this with a silk or rubber belt held tightly between two rollers of different materials. As the belt comes into contact with one roller and then separates, it picks up charge from the roller. When it comes into contact with the other roller and then separates, it deposits the charge. It thus makes one roller positively charged and the other negatively charged. The electrons are supplied to one roller and extracted from the other by metal brushes, connected to the large conductors. This creates a voltage between the conductors, due to their distance. "(Should we really be going into so much detail about it though?... I guess as an example of a very basic electric circuit, it is a good example. Maybe we could just link to Wikipedia triboelectric effect? That's the thing I don't like about wikibooks. You have to provide all of the background instead of just directing them to it.)"

Electronics/Dielectrics. A dielectric is the material between two charges. It impeds the flow of charge until a certain voltage threshold at which point charge starts flowing. "When I first learned about capacitors I thought that this was the way capacitors were used. They charge up to a strong enough voltage and then the dielectric breaks down and then they charge up again. I thought this was the basis for timer circuits. We need to be careful to separate dielectric breakdown from things that are supposed to happen in electronic circuits."

Authoring Webpages/Adapting a webpage for visual browsers. Introduction. The HyperText Mark-up Language (HTML) allows you to add structure and hyperlinks to a text. This in turn allows a web browser to display that text in a useful manner to the user. The mark-up you use has little or nothing to do with the display of the text: this makes it possible to display the hypertext on a wide array of devices. The web browser has to make a translation between the mark-up you provide and display properties. For instance, a heading can be displayed in bold, large text on a graphical browser, and can be spoken out loudly in a speech browser, et cetera. Generally, manufacturers of web browsers make good choices. As well-trained readers, we know that a bold large text on its own line over a mass of normal sized text is probably its heading. So when a browser renders a heading like that, we tend to recognize it as a heading without even thinking. List items are preceded by list item markers such as numbers or 'bullets', emphasized text is printed in bold or italics, et cetera. Manufacturers of graphical web browser even consider that reading off a screen is hard: paragraphs are divided by empty lines, text is generally printed (by default) using a relatively large font-size, and hyperlinks are clearly marked as such, by using different colors and underlines or surrounding lines (the latter in the case of images). There are two things though, that current web browsers are getting wrong. One is that they display each and every webpage the same; the other that they generally fail at presenting webpages in an aesthetically pleasing way. It can certainly be argued that the browser manufacturers cannot be blamed for these problems. Excercise 4-1. Why can manufacturers of web browsers not be blamed for rendering all webpages alike? Why can they not be blamed for displaying webpages in an aesthetically unpleasing way? You may also argue the reverse, if you like. One of the few clues a visitor might have about the quality of webpages is provided by a concept called website. A website is a collection of webpages that belong together. Because they belong together, they provide a powerful hint to the visitor that information that was promised on one page of the website, may be found on another page of that website. Similarly, if the visitor accepts the voice of a webpage author as an authority, other webpages by that author may provide an interesting target for further visits. There are several ways in which an author can make clear that webpages form part of one overarching website. One of these ways has already been discussed: by using a sensible codice_1 text, authors can show that pages belong together. The title of this page, "Adapting a webpage for visual browsers - Wikibooks", helps to underline this. All pages on the Wikibooks website have a title that ends with "- Wikibooks". If you feel like reading other textbooks, the message is clear: you should be on this site. Another way to suggest a relation between webpages, that is, to suggest a website, is in the strategic use of the Uniform Resource Locator (URL), the address of a webpage that is often displayed in a webbrowser's address bar. You can use similar addresses for similar webpages. The third way to suggest a relation between webpages is probably the most useful, though: you can use a uniform visual style to indicate to visitors on which website they are. A bit of history. Originally, HTML was a mish-mash of graphical display mark-up and structural mark-up. This sounds perhaps worse than it was: only a few elements were reserved for visual presentation, such as codice_2 for bold text, and codice_3 for italicized text. Further, codice_4 allowed an author to display text using 'plain text formatting' (as discussed before), codice_5 forced a graphical browser to start displaying following text on a new line, and codice_6 let you display an image at a certain point in the text. All this wasn't so bad, really. It did not introduce a huge drop in accessibility, but allowed authors of pages for graphical clients (more or less the default since the beginning of the web) to 'ponce up' their pages and make them more attractive to visitors. However, by allowing display-only elements in an otherwise display-independent language, Tim Berners-Lee opened the door for abuse by browser manufacturers. The possibility to create a visually pleasing web page made the web a more attractive place to be. Just like the DTP revolution made everybody think that they were graphical designers (mistakenly, most of the time), so did the graphical web open the door for letting form rule over content. When the web became more popular with users, it of course also became more popular with businesses. Actual businesses were founded to produce web browsers, something almost unheard of before. The most successful of these was Netscape. Browser wars. Netscape quickly recognised that the instant appeal of pretty webpages was just as strong a selling-point for a web browser, if not stronger, than the other, more solid appeals, such as the promise of becoming one's own publisher, or being able to traverse a associative landscape of ideas. However, Netscape only had control over the browser, not over the web itself, or its underlying HTM Language. So what Netscape did was let its web browser, Navigator, recognise an extended version of HTML. Frames, for instance, are a Netscape invention, as is JavaScript and the codice_7 element. Using this strategy, Netscape Navigator soon became by far the most used browser on the web. With it the web (and with the web the internet) became a public space, rather than the academic space it had been before. The one software company that had thoroughly missed the internet boat, and that was Microsoft. To the day of writing this, Microsoft still does not understand the internet: they don't see it as a space within which to operate, but rather as a thing to be owned. In the beginning they even tried to replace it with its own 'internet', called MSN; which is of course nonsense, because the internet is not an atomic network. It is a network of network. MSN was to be part of the internet, much to the chagrin of Bill Gates and his people. Meanwhile, Netscape looked further ahead. For them, the internet was a vehicle to gain control over the 'desktop', the metaphor 1980s' operating systems use to describe themselves. If everybody had Netscape Navigator installed, and Netscape Navigator was this universal tool that could be used for everything, from playing games to word processing, it did not really matter which operating system was running beneath the browser. Microsoft was, through its original disdain for the internet and the web, suddenly threatened in the core reason of its existence. It then made a strategic decision that pulled it right into the centre of the internet: it would build its own browser. Microsoft's browser was called Internet Explorer, and it started the so-called browser wars. Microsoft started playing Netscape's game. It first supported almost all of Netscape's HTML extensions, added a few of its own, and changed the behaviour of some elements in a way that made them work slightly better. The web was under a threat of balkanization: the dividing in parts so small, that dealing with the web, and especially authoring web pages, was to become an ordeal that would almost prove too much for authors and surfers alike. Authors suddenly had to decided which browsers to support, or if they should support specific browsers at all. The possibility to treat the web as a purely visual medium planted the misguided thought in a lot of heads that web pages should look alike everywhere. This was when the World Wide Web Consortium stepped in, and started rallying for stricter standards, and a division of document structure and document lay-out. The smartest move of the W3C was to get the browser manufacturers on board. With Microsoft and Netscape both having a direct say about how next versions of HTML would look like, they became stakeholders with an interest in creating a useful web language. The W3C introduced a programming language for suggesting certain lay-outs to a web browser, called CSS. That abbreviation stands for Cascading Style Sheets. Stylesheets had a couple of things going for them. Most importantly, they promised a 'code once, view everywhere' approach to web authoring. This was good for the W3C, because that approach was what HTML was about in the first place. And it was good for the authors, because now they did not have to learn about the quirks of every web browser. Other advantages of CSS were from the start: Future. It could be argued that allowing a few graphical display elements in a hypertext language was a good move. It popularized the web, and with it the internet. It introduced a great deal of authors to the web. However, today many of these authors see the web as a visual medium rather than the (hyper)textual medium it really is. Weaning those people from their wrong notion may prove an impossible task. Websites and home pages. Websites. As we saw before, webpages form part of the web, because they link to other pages or because other pages link to them. In other words, web pages do not live in a vacuum. When you receive a folder in a letter box for a pizza delivery service for instance, there is no context. You do not generally receive copies of competing services simultaneously, or a map that shows you where the delivery service is located, or an encyclopaedia that tells you about the history of pizzas. A webpage does come with such a context. For instance, if you visit a webpage that lets you order pizzas, you probably visited a page before that which let you choose from several pizza delivery services. Also, the order page may link to a map, or to an article about the history of pizzas. Even if it does not do so, the web browser may provide additional functionality to you. For instance, if you selected the name of a pizza in the Firefox web browser, then right-clicked, a menu item would appear that would let you search Google for the selected text in a new tab in the background. Although all webpages have a such a context, only the ones provide by you are under your control. When you group a number of related webpages, such a grouping is called a website. Webpages may form a website for any number of reasons; because they belong to the same subject, because they are hosted on the same server, because they live on the same domain, or because the are created by a single person or organisation. For instance, the collection of webpages at http://www.nasa.gov form the website of the US national space agency NASA. They may not all be served by the same server, and they are not created all by the same person, but they live on the same domain, try to be the voice of NASA on the web and deal with topics that are all related to NASA. A person or an organisation may of course have multiple websites; and what is called a website in this respect is not really important. Websites are often characterised by All these clues tell you after going from one page to another that you are either on the same website, or have left for another website. Most web authors get these clues at least partially right, and most web page visitors are able to tell which website they are on. An interesting feature of a webpage is then that it contains information that is at the core of what the author wanted to say with that particular webpage, but that it also contains information to the visitor that they are on a particular website, where they are on that site, where they can go, etc. Most webpages will contain a lot of information that are pertinent to the subject of that page, and a little information about the website itself. The exceptions are called home pages; home pages are about the website they are part of. Home pages. The "homepage" is the main page of a website. It is used as a central hub for the rest of the site. This is the file that is displayed when going to a web address that doesn't specify a document, e.g. codice_8. It is typically named codice_9, codice_10, codice_11, or codice_12. A homepage has a number of functions and a number of rules and heuristics that it should adhere to. The functions of a homepage are to: Let's review these. Navigational aid. When a visitor follows a information scent to your web page, the web page may or may not fulfill the information wish of the visitor. In the latter case, a visitor will either want to track back, or follow further scents. As we have noted, the indication that a webpage is part of a website provides a powerful hint to the visitor that the same voice that wrote the webpage has written other webpages on possibly similar subjects. A visitor who does not succeed at the current webpage, or who has succeeded, but now has changed goals, may want to further explore your websites. For example, if you collect jokes, and a visitor thinks the first of your jokes is funny, they may wish to read more of your jokes. Links on the webpage to related webpages can be very useful; but sometimes these links are lacking, or they are worded in a way that does not help your visitor, or they do not make clear that they lead to parts of the same website; and often a webpage does not show the intentions an author has with a website, the freshness of a website et cetera. A homepage should provide this sort of information, or at least leads to it. A homepage of a larger should also provide alternative ways of navigation. This could be done through: A common way to link to a homepage is to use the logo or site name as a hyperlink, or to link to it from a breadcrumb trail. Site style. In order to find out that you are on a website, you need to have visited at least two webpages on that site: one to establish that a certain style was used, another to verify that style. Any webpage on a site can be used for that second function, but the homepage must always be usable in this way. Site meaning. The homepage of a website should always make clear what the website is about, just like a webpage should always make clear what that webpage is about. News. In lieu of other designated places where a visitor can review whether pages of a website have been added or updated, the homepage should be regularly updated to indicate that a website is still "alive". Popular ways to do this are to show leads to recent news item, to regularly change highlighted items, or to regularly make simple changes to the lay-out of the homepage. Other ways are to introduce seasonal elements into a homepage. For instance, you could use your homepage to wish visitors happy holidays. Other designated places that allow you to assess the liveliness of a website are the news page, or for example the Recent Changes page at Wikipedia. You can use these pages instead of the homepage for situational feedback, as long as it is clear to visitors that they need to look somewhere else, for instance because you clearly link to a news section. Location. In general, you will wish to let the world know about specific webpages, instead of your website. However, there are instances when the latter is desirable. For those instances it would be useful if the homepage were easy to find and to get to. One way to achieve is, is to locate the homepage at the shortest possible URL. If your webpages are at http://www.example.com/~wily/friends.html, http://www.example.com/~wily/album.html, and http://www.example.com/~wily/contact.html, your homepage should be at http://www.example.com/~wily/. (When faced with a request for a directory instead of a file, a webserver will generally start looking for files with a certain name, such as index.html, index.php, welcome.html etc. This behaviour differs from webserver to webserver, but index.html is usually a pretty safe file name to give to your homepage.) Also, when you forget to link to a homepage, or when the link cannot easily be found, visitors will apply a trick called the directory traversal attack. Despite its goulish name and the fact that it is a crime in the UK, this is a perfectly moral and fine thing to do. It works by guessing the parts of a URL that are superfluous for the homepage address. For example, if you are at http://www.example.com/~wily/friends.html, removing "friends.html" or "~wily/friends.html" from the full address in the address bar of your web browser, may lead to the homepage of this webpage's site. Recapitulation. In essence, a homepage provides situational knowledge about a website. It should show a visitor what the site is about, what its main themes are, how you can get there, how fresh a site is, etc. Cascading Style Sheets. Stylesheets and the style element. CSS (Cascading Style Sheets) is a language used to "style" markup, such as HTML. CSS is a series of rules, and each rule has three parts: a "selector", a "property", and a "value". a { color: red; font-style: italic; In this example, "a" is the selector. It selects all anchors ("links") in the document. Each rule that affects anchors is enclosed in the brackets ({}) following the selector. Here, the two properties are "color" and "font-style". "color" is used to set the color of the text, and "font-style" is used to set the variety of the font. Anchors will appear as red, italicized text. There are three methods of adding CSS to a web page, but the third is considered the best for most purposes. 1) As a tag attribute. The name of the attribute is "style". codice_13 This is considered the worst way of adding styles, in most cases. The reason for this is that it's very hard to maintain. If you wanted to change the color of the anchors on your site, you'd have to find and replace the attribute in every "a" element in every web page, which could occur hundreds or even thousands of times, depending on the size of the website. 2) As a style element. This element is placed inside the codice_14 tag. <stlye> a { color: red; font-style: italic; </style> This is considered the second worst way of adding styles, in most cases. While not nearly as difficult to maintain as the first method, you'd still have to synchronize all of the style elements across all documents, and as your site grows, this would become more and more awkward. 3) As a linked document. The document ends in .css (e.g. "styles.css"). This is considered the best way, in most cases, because it is the easiest to maintain. To link your web page to your stylesheet, add this in your head tag: <link rel="stylesheet" href="styles.css" /> Now, if you want to change any part of the anchors' appearance, you'd only have to edit one file, and it will affect all pages on the website that have that link in the head. Typography on the web. Empty lines separating block level elements--line width----losing control--practice: font family--practice: italics, bold, font-size--practice: line-width--practice: line-height. Colours: dangerous and beautiful. Easy branding with colour--danger: colour blindness--danger: browser settings--solution: never use just colours--solution: always define all colours--practice: applying colours to text and links--practice: applying colour to backgrounds--practice: applying coloured borders. Preparing for print. The print style sheet. Questions and Exercises. Answers. For answers, see Answers to Questions and Exercises. 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Electronics/Permission. Wikibooks GCSE Science. User talk:69.157.40.94 From Wikibooks, the free textbook project. Hi and welcome to wikibooks. I noticed you started a book on electronics and just wanted to point out that i have already written a very basic introduction to electronics in the GCSE science book. Take a look at Circuits (GCSE science) and the pages that follow - you may find it easier to visit Electricity (GCSE science) and navigate from there. Feel free to copy anything you want including any of the diagrams. (click on a diagram and it will take you to a description page titles image:image_name.png, copy the title of the description page into the page you want and put double square brackets around it like this . Preview and you should see the image appear on the page. Theresa knott 12:54, 17 Mar 2004 (UTC) Wikipedia. Do we need any permission to use content from Wikipedia? No. Anything from Wikipedia is copyleft, and can be used as long as the same freedoms are provided. However, you must comply with the site license which includes providing attribution. This is usually done by requesting a history import at . I'm guessing this includes images. No it does not include images. Images hosted on Wikipedia may be subject to copyright and used under fair use criteria that do not necessarily apply here. Electronics textbook. Anyone want to look at this license to see if we can use the electronics book? Klunky. "The web page is written in HTML and Javascript. Because of that the source code is wide open and available to all. Radio amateurs are encourged to take the code, modify it, and create something even more useful to the ham community. The way the web page operates does not have to be limited to schematics. It could be extended to other uses such as printed circuit board layout." "These things imply that it's ideal for amateur radio. It can be shared all over the world, which is where hams are. It can be used with diverse equipment, which hams have. It it can be modified, which hams love to do. And hams need simple, albiet clunky, schematics. So, Klunky is my gift to amateur radio. If some ham, somewhere, can get some use out of it, he's welcome to it." I am sure he wouldn't mind if we use it. He says the code is open. I have written him to ask specifically, but he hasn't responded.

Electronics/Magnetic Field. Magnetic Field: B. A magnetic field is generated due to: formula_1 Permeability: μ. If χ is positive the magnetic field increase and the material is paramagnetic and if χ is negative the magnetic field decreases and the material is diamagnetic. Paramagnetism. Paramagnetism: (with magnetism) tendency of a magnetic dipoles to align itself with an external magnetic field due to the effects of spin. By aligning themselves the dipoles strengthen the external magnetic field. (relative magnetic permeability > 1, magnetic susceptibility 10-3 to 10-5) When the magnetic field is removed thermal motion will quickly disrupt magnetic alignment. Paramagnetism acts on each dipole of an atom separately with no interactions between the dipoles. Behaves like a bar magnet when subject to an external magnetic field. formula_4 M magnetisation, B magnetic flux, T is absolute temperature (Kelvin), and C is a material-specific Curie constant. Increasing the magnetic field makes paramagnetic materials more magnetic, decreasing the temperature reverses this. Curie's is accurate when saturation is low and stops working as the material becomes saturated. Ferromagnetism. Permanent Magnets. Ferromagnetic materials make up permanent magnets. To understand how a permanent magnet is made and destroyed it is necessary to look at the microscopic structure of a magnetic material such as iron. A piece of iron is made up of lots and lots of tiny little magnets called domains. In unmagnetised iron these domains have their north poles pointing all over the place in a random way, but when the iron is put in a strong magnetic field the domains line up with the field. The iron becomes magnetised. Once it has become magnetised, the magnet's own field is strong enough to keep itself magnetised even after it is removed from the external field. "Picture Hysteris Curve" Hysteris curve. heat destroys magnetism if the material is heated above its "Curie" point. and becomes more aligned. When the field is removed the magnet retains some of the field and is magnetized. Moving Charge. When a negatively charged particle such as an electron moves through space or a wire it generates a magnetic field that circles counterclockwise in a plane perpendicular to it. This is part of Ampere's law. For a positively charged particle the magnetic field is clockwise. The faster the charged particle moves the larger the magnetic field. Inductance. Self inductance This is what happens when you turn on DC in an inductor. All remains calm until you try to turn the thing off For large coils the EMF can be huge, sparks can fly! Isn't really an issue for DC though except when turning on and off. But for AC the current is varying all the time, so things never settle into a steady state. The higher the frequency, the faster things are changing and so the greater the inductance. I hope all this is reasonably clear now. Theresa knott 13:49, 30 Apr 2004 (UTC) Electromagnet. A magnetic hose that acts like a bar magnet. As the current increases the magnet becomes more repulsive to other magnets. Used as a switch for magnetic materials and to pick up magnetic material such as cars at scrapyards. To make a good electromagnet we use a 'soft iron' core. Soft iron has the property of being easy to magnetise. It also loses its magnetism very easily once the electric current is turned off. Steel can also be used as the core of an electromagnet, but steel does not lose its magnetism once the electric current is turned off. It stays magnetised. When a current is switched on, or off, it creates a changing magnetic field, which in turn creates voltages that oppose the current. This is the basis for inductance. By coiling the wire up we can make a solenoid. The magnetic field of a solenoid looks like the field of an ordinary bar magnet. A solenoíd makes a pretty weak bar magnet on its own but if a piece of iron is put inside the solenoid the field becomes much much stronger. Try the following experiment: You will find that the solenoid with the iron core is very much stronger than the one with the air core. Q1) Why is it important to use plastic coated wire? Mutual Inductance. Mutual inductance is just like inductance except that the magnetic field from the current in one solenoid is creating voltage in a "different" solenoid, as well as in itself. Toroid. "Picture of a toroid" A toroid is a solenoid that wraps around on itself. toroid: A donut shaped inductor where the magnetic fields stay contained within it. Has an air gap. A Toroid is a doughnut-shaped object enclosed by a torus. It has a ring-shaped surface generated by rotating a circle around an axis that does not intersect the circle. Usually, a toroid is composed of a coil of insulated wire in a donut shape (and is usually composed of iron or similar metals). Toroids are inductors in circuits. They are primarily used in low frequencies transception for the inductance properties needed. Toroids possess greater inductance properties and carry greater current than similarly constructed solenoids. Toroidal coils reduce resistance, due to diameter and number. In a toroid, magnetic fluxes are confined in the toroid's core. a toroid magnifies current. Magnetic Fields and Electric Fields. Notice the similarity between magnetism and electricity. Electricity attracts neutral material (recall a rubbed balloon picking up small bits of paper). Electricity comes in two forms: positive and negative. Like charges repel unlike, charges attract. Physicists long ago realised that these similarities were important. It is now realised that electricity and magnetism are closely related forces. They are really two different aspects of the same force! We call that force electromagnetism. This next section of the electricity module is probably the hardest, but it is also the most interesting. Work out the problems in the text as you go along and you "will" be able to cope. A "changing" electric field creates a magnetic field. Conversely, a "changing" magnetic field creates an electric field. used to create photon / radio waves This means that electric currents produce magnetic fields and changing or moving magnetic fields produce electric currents. This module looks at this property of electric and magnetic fields in detail. You will learn about the magnetic fields generated by a wire and a solenoid, then go on to look at motors, generators, and transformers. Any charged particle in a magnetic field feels a force perpendicular to both the field and its own velocity. So far we have looked at the effect of putting a current in a magnetic field and seeing that there is a force on the wire that carries the current. On this page we will look at a related effect known as "induction". Talk about straight magnetic fields and electrons moving in circular path. Talk about how magnetic fields alter the paths of charged particles. For a more math oriented explanation of Electric and Magnetic Fields consult Electrodynamics. Theoretical explanation of what's going on (Advanced). So now we know what happens we have to come up with an explanation. Look at the diagram below. In this diagram the green line represents a wire that is going into the screen or page. It is at right angles to the magnetic field between the two bar magnets. The wire is moved downwards so that it "cuts" the field lines. This induces an e.m.f. (voltage) along the wire. It is this e.m.f. that causes a current to flow if a complete circuit is made. So now we know why you get a current. It's caused by the induced voltage, but why do you get a voltage? This is a close up of the wire seen end-on. The blue circles represent where the magnetic field lines stick out of the screen or page. They are coming straight out of the screen. The wire is made of a metal and so has many free electrons. As the wire moves downwards each individual electron travels downwards too. They effectively form loads of little current flowing down. From the motor effect we know that if a current flows in a magnetic field there will be a force. In this case the force pushes the electrons to the right. It is the electrons all moving up to the right end of the wire that causes the voltage difference. So, the motor effect can induce an e.m.f. This is one of two ways magnetic fields can induce an e.m.f Q3) Imagine more than one loop of wire cutting through the magnetic field.Each loop of wire has its own e.m.f. induced on it. If the wires are joined up, it's as if they are little cells in series. Fill in the blanks " As the number of windings on the solenoid "increases" the value of the induced voltage ______________. When a current flows down a wire it creates a magnetic field. To see this place a small plotting compass near the wire and turn the current on.The compass needle will deflect. The shape of the magnetic field around a wire is circular. look at the diagram on the right. The wire is coming straight out of the screen so you only see it's cross section (the red circle) The plotting compasses show how the field wraps around the wire. Making the field stronger. To make the field stronger we can: Reversing the field. To make the north pole and south pole swap positions we can Talking of north and south poles, there is a little trick to find out which end will be the north pole and which end will be the south. Look down the solenoid and work out which way the current is flowing. Remember current flows from the positive to the negative terminals on the power pack. If the current is flowing clockwise the end will be a south pole. If it looks like it's going anticlockwise the end will be a north pole. The easy way to remember this is to put arrows on the end of a capital N and S like this.

Electronics/Spin. Electrons posses a quantum mechanical property called spin. It is responsible for magnetism. It works when charge is moving. When charge does not flow there is no magnetic field. Magnetic fields work in a direction perpendicular to the flow of charge. Sure, it's the easiest way to understand magnetism. We don't need to talk about the formulas associated with spin, we just need to give an overview of how it works.

Electronics/Charge. Two atoms are walking down the street. The first atom says to the second atom "I think I lost an electron!" The second says "Are you sure?" To which the first states "I'm positive!" Balance of Charge. Atoms, the smallest particles of matter, are made of protons, electrons, and neutrons. Protons have a positive charge, Electrons have a negative charge that cancels the proton's positive charge. Neutrons are particles that are similar to a proton but have a neutral charge. There are no differences between positive and negative charges except that particles with the same charge repel each other and particles with opposite charge attract each other. If a solitary positive proton and negative electron are placed near each other they will come together to form a hydrogen atom. This repulsion and attraction (force between stationary charged particles) is known as Electrostatic Force and extends theoretically to infinity, but is diluted with distance. Both atoms and the universe have a neutral charge overall and come with the same number of protons and electrons. When an atom has one or more missing electrons it is left with a positive charge, and when an atom has at least one extra electron it has a negative charge. Having a positive or a negative charge makes an atom an ion. Atoms only gain and lose protons and neutrons through fusion, fission, and being radioactive. Although atoms are made of many particles and objects are made of many atoms, they behave similar to charged particles in terms of how they repel and attract. In an atom the protons and neutrons combine to form a tightly bound nucleus. This nucleus is surrounded by a vast cloud of electrons circling it at great distance held near the protons by electromagnetic attraction. The cloud exists as a series of overlapping shells / bands in which the inner valence bands are filled with electrons and are tightly bound to the atom. The outer conduction bands contain no electrons except those that have accelerated to the conduction bands by gaining energy. With enough energy an electron will escape an atom. When an electron in the conduction band decelerates and falls to another conduction band or the valence band emitting an photon, is known as the photoelectric effect. Electronic devices are based on the idea of exploiting the differences between conductors, insulators, and semiconductors. Conductors. In a conductor the electrons of an object are free to move. Due to their repulsive nature the valence electrons leave the center of the object and spread out evenly across its surface in order to be as far apart as possible. This cavity of empty space is known as a Faraday Cage and stops radiation, such as charge, radio waves, and EMPs (Electro-Magnetic Pulses) from entering and leaving the object. If there are holes in the Faraday Cage than radiation can pass. One of the interesting things to do with conductors is the transfer of charge between metal spheres. At the start the spheres are neutral. The first step involves putting sphere 1 next to but not touching sphere 2. This causes all the electrons in sphere 2 to travel away from sphere 1 to the far end of sphere 2. So sphere 2 now has a negative end filled with electrons and a positive end lacking electrons. Next sphere 2 is grounded by contact with a conductor connected with the earth and the earth takes its electrons leaving sphere 2 with a positive charge. The positive charge (absence of electrons) spreads evenly across the surface due to its lack of electrons. Insulators. In an insulator the electrons of an object are stuck. This allows charge to build up on the surface of the object by way of the triboelectric effect. The triboelectric effect (rubbing electricity effect) involves the exchange of electrons when two different insulators such as glass, hard rubber, amber, or even the seat of one's pants, come into contact. The polarity and strength of the charges produced differ according to the material composition and its surface smoothness. For example, glass rubbed with silk will build up a charge, as will hard rubber rubbed with fur. The effect is greatly enhanced by rubbing materials together. Because the material being rubbed is now charged, contact with an uncharged object or an object with the opposite charge may cause a discharge of the built-up static electricity by way of a spark. A person simply walking across a carpet may build up enough charge to cause a spark to travel over a centimeter. The spark is powerful enough to attract dust particles to cloth, destroy electrical equipment, ignite gas fumes, and create lightning. In extreme cases the spark can destroy factories that deal with gunpowder and explosives. The best way to remove static electricity is by discharging it through grounding. Humid air will also slowly discharge static electricity. This is one reason why cells and capacitors lose charge over time. Quantity of Charge. Protons and electrons have opposite but equal charge. Because in almost all cases, the charge on protons or electrons is the smallest amount of charge commonly discussed, the quantity of charge of one proton is considered one positive elementary charge and the charge of one electron is one negative elementary charge. Because atoms and such particles are so small, and charge in amounts of multi-trillions of elementary charges are usually discussed, a much larger unit of charge is typically used. The Coulomb is a unit of charge, which can be expressed as a positive or negative number, which is equal to approximately 6.2415×1018 elementary charges. Accordingly, an elementary charge is equal to approximately 1.602 ×10-19 coulombs. The commonly used abbreviation for the Coulomb is a capital C. The SI definition of a Coulomb is the quantity of charge which passes a point over a period of 1 second ( s ) when a current of 1 Ampere ( A ) flows past that point, i. e., C = A s or A = C/s. An is one of the fundamental units in physics from which various other units are defined, such as the coulomb.

Electronics/Frequency Spectrum. This page covers the spectrum of electromagnetic waves. Sound. Sound is heard when there are two that things collide with each other . Sound wave is a longitude wave . Sound waves are affected by temperature and atmosphere's pressure In water sound waves travel faster due to particles being more compressed. In Air, Sound wave pressure creates rare-fraction and rare-compression columns of air dB scale. DB Scale is used to measured the pressure of a Sound Infrasound: < 20 Hz. Sounds of frequency below 20 Hz is called InfraSound can not be heard by human's ears sound travels the earth so it can be used to detect nuclear blasts and allows whales to communicate for 10,000 miles (sound channel). can be greater than >120 dB is intensity but not heard by people. associated with avalanches, earthquakes, volcanoes, meteors, the earth groaning, rock concerts, and pipe organs. Intensifies emotions. Infrasound has some really warped effects on society. sound that travels the earth often >120 dB is intensity. If heard would be a large rumble as loud as a jet engine or gunfire. At rock concerts they play infrasound. This intensifies emotions. This is one possible reason why people play loud music. Attraction of pipe organs. In addition the earth groans. "Anyone know the mythological significance of the cry of the earth?" used by whales to communicate as the waves can travel 10,000 miles. Infrasound appears to heighten emotions and make people nervous and thus has an association with the supernatural such as haunted houses. Since people are emotionally attached to the deceased and go to where they died, when you combine it with infrasound which is associated with chills, is this an explanation for ghosts? One of the stranger uses of infrasound involves combining an ancient chamber that resonates infrasound with mind altering drugs to induce visions and out of body experiences. Audible Sound: 20-20,000 Hz. Audible Sound are sound that can be perceived by human's ears . Audible Sounds are sounds in the range of frequencies 20-20,000 Hz Human's Voice Computer Sound formats Animal Sound Music Telephone Ultrasound: > 20 kHz. Sounds above 20 kHz are called Ultra Sound . Unable to be heard by people, but can be heard by dogs and dolphins. Light. photons are emitted by a black body. Talk about energy of photons and how far they penetrate. regulated by ITU In the US, the allocation of these bands is managed by the FCC Federal Standard 1037C and MIL-STD-188. The VLA telescope array can resolve a &lambda of 8,xxx Hz Candela. Unit of intensity Radio Waves < 3 GHz. (HF) High Frequency: 3-30 MHz. * Sunlight/darkness at site of transmission and reception * Transmitter/receiver proximity to terminator * Season * Sunspot cycle * Solar activity * Polar aurora * Maximum usable frequency * Lowest usable frequency * Frequency of operation within the HF range * The distance from the transmitter to the target receiver * Time of day. During the day, higher shortwave frequencies (> 10 MHz) can travel longer distances than lower; at night, this : property is reversed. * Season of the year. * Solar conditions, including the number of sunspots, solar flares, and overall solar activity. Solar flares can prevent the : ionosphere from reflecting or refracting radio waves. * Domestic broadcasting in countries with a widely dispersed population with few longwave, mediumwave or FM stations serving them * International broadcasting to foreign audiences * Utility stations transmitting messages not intended for a general public, such as aircraft flying between continents, encoded or ciphered diplomatic messages, weather reporting, or ships at sea * Amateur radio operators * Time signal stations 5,900-5,950 kHz 7,300-7,350 kHz (7,100-7,350 kHz outside of the Americas) 9,400-9,500 kHz 11,600-11,650 kHz 12,050-12,100 kHz 13,570-13,600 kHz 13,800-13,870 kHz 15,600-15,800 kHz 17,480-17,550 kHz 18,900-19,020 kHz Table of contents [showhide] 1 International broadcasting 2 Amateur radio 3 Shortwave listening 4 Numbers Stations 5 Shortwave radio: the future 1 Shortwave Broadcasts and Music International broadcasting See International broadcasting for details on the history and practice of broadcasting to foreign audiences. Amateur radio The privilege of operating shortwave radio transmitters for non-commercial purposes is open to licensed amateurs. In the US, they are licensed by the Federal Communications Commission (FCC). U.S. citizens do not need licenses to own or operate shortwave receivers. Recently the FCC has added an amateur radio license which requires no knowledge of Morse code, making it easier for beginners to get involved; however, a working knowledge of Morse code is required to operate on shortwave bands. Amateur radio operators have made numerous technical advancements in the field of radio and make themselves available to transmit emergency communications when normal communications channels fail. Some amateurs practice operating off the power grid so as to be prepared for power loss. Shortwave listening Many hobbyists listen to shortwave broadcasters without operating transmitters. In many cases, the goal is to obtain as many stations from as many countries as possible (DXing); others listen to specialized shortwave utility, or "ute", transmissions such as maritime, naval, aviation, or military signals. Others focus on intelligence signals. Shortwave listeners, or SWLs, can obtain "QSL" cards from broadcasters or utility stations as trophies of the hobby. Numbers Stations Numbers stations are shortwave radio stations of uncertain origin that broadcast streams of numbers, words, or phonetic sounds. Although officially there is no indication of their origin, radio hobbyists have determined that many of them are used by intelligence services as one-way communication to agents in other countries. Shortwave radio: the future The development of direct broadcasts from satellites has reduced the demand for shortwave receivers, but there are still a great number of shortwave broadcasters. A new digital radio technology, Digital Radio Mondiale, is expected to improve the quality of shortwave audio from very poor to standards comparable to the FM broadcast band. The future of shortwave radio is threatened by the uprise of power line communication (PLC), where a data stream is transmitted over unshielded power lines. As the frequencies used overlap with the shortwave bands severe distortions make listening to shortwave radio near power lines difficult or impossible. Shortwave Broadcasts and Music Some musicians have been attracted to the unique aural qualities of shortwave radio. John Cage employed shortwave radios as live instruments in a number of pieces, and other musicians have sampled broadcasts, used tape loops of broadcasts, or drawn inspiration from the unusual sounds on some frequencies. Karlheinz Stockhausen used shortwave radio in works including Telemusik (1966), Hymnen (1966–67) and Spiral (1968), and Holger Czukay, Pat Metheny, Aphex Twin, Meat Beat Manifesto, and Wilco have also used broadcasts. (VHF) Very High Frequency: 30-300 MHz/10-1 m. FM been around since 1940 42-50 MHz Pre WWII allocation of FM 1940 88-106 MHz Post WWII allocation of FM 1945 FM stereo technology Mazer = microwave (UHF) Ultra High Frequency: 300-3000 MHz/100-10 cm. UHF and VHF are the most common frequency bands for television. UHF frequencies have higher attenuation from atmospheric moisture and benefit less from 'bounce', or the reflection of signals off the ionosphere back to earth, when compared to VHF frequencies. The frequencies of 300-3000 MHz are always at least an order of magnitude above the MUF (Maximum Usable Frequency). The MUF for most of the earth is generally between 25-35 MHz. Higher frequencies also benefit less from ground mode transmission. However, the short wavelengths of UHF frequencies allow compact receiving antennas with narrow elements; many people consider them less ugly than VHF-receiving models In the United States, UHF stations (broadcast channels above 13) originally gained a reputation for being more locally owned, less polished, less professional, less popular, and for having a weaker signal than their VHF counterparts (channels 2-13). The movie UHF, starring Weird Al Yankovic, parodies this phenomenon. Recently, with the emergence of eight major broadcast television networks, that notion has changed as bigger and bigger media companies seek a bigger slice of the television pie. Many Fox, UPN, WB, and Pax network affiliates broadcast in the UHF band. As cable television, digital television, and DSS have penetrated the television market, the distinction between VHF and UHF stations has dissipated. In Australia, UHF was first anticipated in the mid 1970s with channels 28 to 69. The first UHF TV broadcasts in Australia were operated by Special Broadcasting Service (SBS) on channel 28 in Sydney and Melbourne starting in 1980, and translator stations for the Australian Broadcasting Corporation (ABC). The UHF band is now used extensively as ABC, SBS, commercial and community (public access) television services have expanded particularly through regional areas. Ultra High Frequency A band reserved for television. Microwave 3-300 GHz/100-1 mm. (SHF) Super High Frequency: 3-30 GHz/10-1 cm. Microwave Note: above 300 GHz, the absorption of electromagnetic radiation by Earth's atmosphere is so great that the atmosphere is effectively opaque to higher frequencies of electromagnetic radiation, until the atmosphere becomes transparent again in the so-called infrared and optical window frequency ranges. Uses A microwave oven uses a magnetron microwave generator to produce microwaves at a frequency of approximately 2.4 GHz for the purpose of cooking food. Microwaves cook food by causing molecules of water and other compounds to vibrate. The vibration creates heat which warms the food. Since organic matter is made up primarily of water, food is easily cooked by this method. A maser is a device similar to a laser, except that it works at microwave frequencies. Microwaves are also used in satellite transmissions because this frequency passes easily through the Earth's atmosphere with less interference than higher wavelengths. Radar also uses microwave radiation to detect the range, speed, and other characteristics of remote objects. Wireless LAN communication protocols such as IEEE 802.11 and bluetooth also use microwaves in the 2.4 GHz ISM band, although some variants use a 5 GHz band for communication. Cable TV and Internet access on coax cable as well as broadcast television use some of the lower microwave frequencies. Microwaves can be used to transmit power over long distances, create a solar array to beam power to earth. The microwave spectrum is defined as electromagnetic energy ranging from approximately 300 MHz to 1000 GHz in frequency. Most common applications are within the 1 to 40 GHz range. Microwave Frequency Bands are defined in the table below: L Band 1 to 2 GHz S Band 2 to 4 GHz C Band 4 to 8 GHz X Band 8 to 12 GHz Ku Band 12 to 18 GHz K Band 18 to 26 GHz Ka Band 26 to 40 GHz Q Band 30 to 50 GHz U Band 40 to 60 GHz V Band 46 to 56 GHz W Band 56 to 100 GHz For some of the history in the development of electromagnetic theory applicable to modern microwave applications see the following figures: * Michael Faraday. * James Clerk Maxwell. * Heinrich Hertz. * Nikola Tesla. * Guglielmo Marconi. * Samuel Morse. * Sir William Thomson, later Lord Kelvin. * Oliver Heaviside. * Lord Rayleigh. * Oliver Lodge. Specific significant areas of research and work developing microwaves and their applications: Specific work on microwaves Work carried out by Area of work Barkhausen and Kurz Positive grid oscillators Hull Smooth bore magnetron Varian Brothers Velocity modulated electron beam → klystron tube Randall and Boot Cavity magnetron See also: * Home appliances. * Microwave auditory effect. * Radio. * Optics. The Cosmic Microwave Background Radiation is a form of electromagnetic radiation that fills the whole of the universe. It has the characteristics of black-body radiation at a temperature of 2.726 kelvins. It has a frequency in the microwave range. CMR and the Big Bang This radiation is regarded as the best available evidence of the Big Bang (BB) theory and its discovery in the mid-1960s marked the end of general interest for alternatives such as the steady state theory. The CMR gives a snapshot of the Universe when, according to standard cosmology, the temperature dropped enough to allow electrons and protons to form hydrogen atoms, thus making the universe transparent to radiation. When it originated some 300,000 years after the Big Bang—this point in time is generally known as the "last scattering surface"—the temperature of the Universe was about 6000 K. Since then it has dropped because of the expansion of the Universe, which cools radiation inversely proportional to the fourth power of the Universe's scale length. Features One of the microwave background's most salient features is a high degree of isotropy. There are some anisotropies, the most pronounced of which is the dipole anisotropy at a level of about 10-4 at a scale of 180 degrees of arc. It is due to the motion of the observer against the CBR, which is some 700 km/s for the Earth. Variations due to external physics also exist; the Sunyaev-Zel'dovich Effect is one of the major factors here, in which an cloud of high energy electrons scatters the radiation transferring some energy to the CMB photons. Even more interesting are anisotropies at a level of roughly 1/100000 and on a scale of a few arc minutes. Those very small variations correspond to the density fluctuations at the last scattering surface and give valuable information about the seeds for the large scale structures we observe now. These density fluctations arise because different parts of the universe are not in contact with each other. In addition, the Sachs-Wolfe effect causes photons from the Cosmic microwave background to be gravitationally redshifted. These small-scale variations give observational constraints on the properties of universe, and are therefore one important test for cosmological models. Detection, Prediction and Discovery The CBR was predicted by George Gamow, Ralph Alpher, and Robert Hermann in the 1940s and was accidentally discovered in 1964 by Penzias and Wilson, who received a Nobel Prize for this discovery. The CBR had, however, been detected and its temperature deduced in 1941, seven years before Gamow's prediction. Based on the study of narrow absorption line features in the spectra of stars, the astronomer Andrew McKellar wrote: "It can be calculated that the 'rotational' temperature of interstellar space is 2 K." Because water absorbs microwave radiation, a fact that is used to build microwave ovens, it is rather difficult to observe the CMB with ground-based instruments. CMB research therefore makes increasing use of air and space-borne experiments. Experiments Of these experiments, the Cosmic Background Explorer (COBE) satellite that was flown in 1989-1996 is probably the most famous and which made the first detection of the large scale anisotropies (other than the dipole). In June 2001, NASA launched a second CBR space mission, WMAP, to make detailed measurements of the anisotropies over the full sky. Results from this mission provide a detailed measurement of the angular power spectrum down to degree scales, giving detailed constraints on various cosmological parameters. The results are broadly consistent with those expected from cosmic inflation as well as various other competing theories, and are available in detail at NASA's data center for Cosmic Microwave Background (CMB) [ed. see links below], A third space mission, Planck, is to be launched in 2007. Unlike the previous two space missions, Planck is a collaboration between NASA and ESA (the European Space Agency). CBR and Non-Standard Cosmologies During the mid-1990's, the lack of detection of anisotropies in the CBR led to some interest in nonstandard cosmologies (such as plasma cosmology) mostly as a backup in case detectors failed to find anisotropy in the CBR. The discovery of these anisotropies combined with a large amount of new data coming in has greatly reduced interest in these alternative theories. Some supporters of non-standard cosmology argue that the primodorial background radiation is uniform (which is inconsistent with the big bang) and that the variations in the CBR are due to the Sunyaev-Zel'dovich effect mentioned above (among other effects). (IR) Infared 1 mm-700 nm. Note: above 100 GHz, the absorption of electromagnetic radiation by Earth's atmosphere is so great that the atmosphere is effectively opaque to higher frequencies of electromagnetic radiation, until the atmosphere becomes transparent again in the so-called infrared and optical window frequency ranges. Infra-red means below red, where red is the color with the longest wavelength. (MIR/IIR) mid-IR (5-30 µm) - Infrared radiation is often linked to heat, since objects at room temperature or above will emit radiation mostly concentrated in the mid-infrared band (see black body). However, these terms are not precise, and are used differently in various studies. Uses Infrared is used in night-vision equipment, when there is insufficient visible light to see an object. The radiation is detected and turned into an image on a screen, hotter objects showing up brighter, enabling the police and military to chase targets. "Night vision" Smoke is more transparent to infrared than to visible light, so fire fighters use infrared imaging equipment when working in smoke-filled areas. Fire fighters also use this equipment in wood-frame buildings after a fire has been extinguished to look for hot spots behind the walls, where a fire can break out again. A more common use of IR is in television remote controls. In this case it is used in preference to radio waves because it does not interfere with the television signal. IR data transmission is also employed in short-range communication among computer peripherals and personal digital assistants. These devices usually conform to standards published by IrDA, the Infrared Data Association. Remote controls and IrDA devices use infrared light-emitting diodes (LEDs) to emit infrared radiation which is focused by a plastic lens into a narrow beam. The beam is modulated, i.e. switched on and off, to encode the data. The receiver uses a silicon photodiode to convert the infrared radiation to an electric current. It responds only to the rapidly pulsing signal created by the transmitter, and filters out slowly changing infrared radiation from sunlight, people and other warm objects. The light used in fiber optic communication is typically infrared. heat based "not really". EM waves are constantly being emitted by objects, black-body radiation, but the frequency depends on temperature. for room temperature objects and humans, it is mostly below light in the infrared, but for things like fires and red hot iron, it is in the light spectrum more, and for colder and hotter objects it is at different frequencies. how about this, anything that has at least as much energy as infared will warm things up. Anything that has less will not warm atoms up. Global Warming Carbon Dioxide absorbs certain wavlengths of infrared light. This is thought to cause global warming. Visible. 1 PetaHz 1015 Hz Visible light exists on the tip end of infared and just below it is UV. Basically you can almost get a sunburn from visible light. The optical spectrum (visible light or visible spectrum) is the portion of the electromagnetic spectrum that is visible to the human eye. The optical spectrum is a composite, or mixture, of the various colors. There are no exact bounds to the optical spectrum ; a light-adapted eye typically has a maximum sensitivity of ~555 nm (in the green). Commonly the response of the eye is considered to cover 380 nm to 780 nm although a range of 400 nm to 700 nm range is more common. The eye may, however, have some visual response at even wider wavelength ranges. Wavelengths in the range visible to the eye occupy most of the "optical window", a range of wavelengths that are easily transmitted through the Earth's atmosphere. Note: Ultraviolet and Infrared are often considered to be "light" but are generally not visible to the human eye. Visible Light Visible light is electromagnetic radiation that is made of colors of light that the eye can see. This light has wavelengths that are generally expressed in nanometers. See Rydberg formula. Vision Rods and cones People see red, yellow, green, blue and purple People don't notice the difference because red and purple both have a low intensity. purple is cut off because the atmosphere absorbs it. light bulbs (UV) UltraViolet. Ultra-Violet means beyond violet, where violet is the color with the shortest wavelength. It is colloquially called black light, as it is invisible to the human eye. The sun emits ultraviolet light in the UV-A, UV-B, and UV-C bands Ordinary glass is transparent to UV-A but is opaque to shorter wavelengths. This is why people don't get sunburns in cars. During the first nuclear blast, Feynman was the only person who saw it because he looked through a window. Quartz glass, depending on quality, can be transparent even to vacuum UV wavelengths. Table of contents [showhide] 1 Health effects 2 Astronomy 3 Uses 3.1 Fluorescent lamps 3.2 Disinfecting drinking water 3.3 Analyzing minerals 3.4 Sterilization 3.5 Resolution 3.6 Spectroscopy 3.7 Photolithography 3.8 Other 4 References Health effects In general, UV-A is the least harmful, but can contribute to the aging of skin, DNA damage and possibly skin cancer. It penetrates deeply and does not cause sun burn. Because it does not cause reddenning of the skin (erythema) it can not be measured in the SPF testing. There is no good clinical measurement of the blocking of UVA radiation, but it is important that your sunscreen block both UVA and UVB. High intensities of UV-B light are hazardous to the eyes, and exposure can cause welder's flash (photokeratitis or arc eye). UVA, UV-B and UV-C all can damage collagen fibers and thereby accelerate aging of the skin. Tungsten-Halogen lamps have bulbs made of quartz, not of ordinary glass. Tungsten-Halogen lamps that are not filtered by an additional layer of ordinary glass are a common, useful, and possibly dangerous, source of UV-B light. UV-A light is known as "dark-light" and, because of its longer wavelength, can penetrate most windows. It also penetrates deeper into the skin than UV-B light and is thought to be a prime cause of wrinkles. UV-B light in particular has been linked to skin cancers such as melanoma. The radiation ionizes DNA molecules in skin cells, causing covalent bonds to form between adjacent thymine bases, producing thymidine dimers. Thymidine dimers do not base pair normally, which can cause distortion of the DNA helix, stalled replication, gaps, and misincorporation. These can lead to mutations, which can result in cancerous growths. The mutagenicity of UV radiation can be easily observed in bacteria cultures. This cancer connection is the reason for concern about ozone depletion and the ozone hole. UV-C rays are the strongest, most dangerous type of ultraviolet light. Little attention has been given to UV-C rays in the past since they are normally filtered out by the ozone layer and do not reach the Earth. Thinning of the ozone layer and holes in the ozone layer are causing increased concern about the potential for UVC light exposure, however. As a defense against UV light, the body tans when exposed to moderate (depending on skin type) levels of radiation by releasing the brown pigment melanin. This helps to block UV penetration and prevent damage to the vulnerable skin tissues deeper down. Suntan lotion that partly blocks UV is widely available (often referred to as "sun block" or "sunscreen"). Most of these products contain an "SPF rating" that describes the amount of protection given. This protection applies only to UV-B light. In any case, most dermatologists recommend against prolonged sunbathing. A positive effect of UV light is that it induces the production of vitamin D in the skin. Grant (2002) claims tens of thousands of premature deaths occur in the U.S. annually from cancer due to insufficient UV-B exposures (apparently via vitamin D deficiency). Astronomy In astronomy, very hot objects preferentially emit UV light (see Wien's law). However, the same ozone layer that protects us causes difficulties for astronomers observing from the earth, so most UV observations are made from space. (See UV astronomy, space observatory). Uses UV light has many various uses: Fluorescent lamps Fluorescent lamps produce UV light by the emission of low-pressure mercury gas. A phosphorescent coating on the inside of the tubes absorbs the UV and turns it into visible light. The main mercury emission wavelength is in the UV-C range. Unshielded exposure of the skin or eyes to mercury arc lamps that do not have a conversion phosphor is quite dangerous. The light from a mercury lamp is predominantly at discrete wavelengths. Other practical UV light sources with more continuous emission spectra include xenon arc lamps (commonly used as sunlight simulators), deuterium arc lamps, mercury-xenon arc lamps, metal-halide arc lamps, and tungsten-halogen incandescent lamps. Disinfecting drinking water Ultraviolet light is increasingly being used to disinfect drinking water and in waste water treatment plants. Dr. James R. Bolton discovered that ultraviolet light could treat Cryptosporidium, previously unknown. The findings resulted in two US patents and the use of UV light as a viable method to treat drinking water. Analyzing minerals Ultraviolet lamps are also used in analyzing minerals, gems, and in other detective work including authentication of various collectibles. Materials may look the same under visible light, but fluoresce to different degrees under ultraviolet light; or may fluoresce differently under short wave ultraviolet versus long wave ultra violet. UV fluorescent dyes are used in many applications (for example, biochemistry and forensics). The fluorescent protein Green Fluorescent Protein (GFP) is often used in genetics as a marker. Many substances, proteins for instance, have significant light absorption bands in the ultraviolet that are of use and interest in biochemistry and related fields. UV-capable spectrophotometers are common in such laboratories. Sterilization Ultraviolet lamps are used to sterilize workspaces and tools used in biology laboratories and medical facilities. Since microorganisms can be shielded from ultraviolet light in small cracks and other shaded areas, however, these lamps are used only as a supplement to other sterilization techniques. Resolution Ultraviolet light is used for very fine resolution photolithography, as required for manufacture of semiconductors. Spectroscopy UV light is often used in UV-visible spectroscopy Photolithography UV light is used extensively in the electronics industry in a procedure known as photolithography, where a chemical known as a photoresist is exposed to UV light which has passed through a mask. The light allows chemical reactions to take place in the photoresist, and after development (a step that either removes the exposed or unexposed photoresist), a geometric pattern which is determined by the mask remains on the sample. Further steps may then be taken to "etch" away parts of the sample with no photoresist remaining. It is photolithography which is primarily used to create integrated circuit components and printed circuit boards. Other The onset of vacuum UV, 200 nm, is defined by the fact that ordinary air is opaque below this wavelength. This opacity is due to the strong absorption of light of these wavelengths by oxygen in the air. Pure nitrogen (less than about 10 ppm oxygen) is transparent to wavelengths in the range of about 150-200 nm. This has wide practical significance now that semiconductor manufacturing processes are using wavelengths shorter than 200 nm. By working in oxygen-free gas, the equipment does not have to be built to withstand the pressure differences required to work in a vacuum. Some other scientific instruments, such as Circular Dichroism spectrometers, are also commonly nitrogen purged and operate in this spectral region. It is advisable to use protective eyewear when working with ultraviolet light, especially short wave ultraviolet. Ordinary eyeglasses give some protection. Most plastic lenses give more protection than glass lenses. Some plastic lens materials, such as polycarbonate, block most UV. There are protective treatments available for eyeglass lenses that need it to give better protection. The most important reason that ordinary eyeglasses only give limited protection, however, is that light can reach the eye without going through the lens. Full coverage is important if the risk from exposure is high. Full coverage eye protection is usually recommended for high altitude mountaineering, for instance. Mountaineers are exposed to higher than ordinary levels of UV light, both because there is less atmospheric filtering and because of reflection from snow and ice. Some insects, such as bees, can see into the near ultraviolet, and flowers often have markings visible to such pollinators. UV catastrophe at wavelengths just below visible is where people get sunburns. sun burn rating UV-A UV-B UV-C UV bulbs X-Ray: 30 PetaHz - 60 ExaHz/10 nanometer - 5 picometer. Soft X-ray: (10-0.1 nm) Hard X-ray: (100-5 pm)Low energy gamma rays X-rays are generated by electron collisions. X-rays are primarily used for diagnostic medical imaging and crystallography. Physicist Johann Hittorf observed tubes with energy rays extending from a negative electrode. These rays produced a fluorescence when they hit the glass walls of the tubes. In 1876 the effect was named "cathode rays" by Eugene Goldstein. Later, English physicist William Crookes investigated the effects of energy discharges on rare gases. He constructed what is called the Crookes tube. It is a glass vacuum cylinder, containing electrodes for discharges of a high voltage electric current. He found, when he placed unexposed photographic plates near the tube, that some of them were flawed by shadows, though he did not investigate this effect. In 1892, Heinrich Hertz began experimenting and demonstrated that cathode rays could penetrate very thin metal foil (such as aluminum). Philip Lenard, a student of Heinrich Hertz, further researched this effect. He developed a version of the cathode tube and studied the penetration of X-rays through various materials. Philip Lenard, though, did not realize that he was producing X-rays. In April 1887, Nikola Tesla began to investigate X-rays using his own devices as well as Crookes tubes. Tesla did this by experimenting with high voltages and vacuum tubes. From Nikola Tesla's technical publications, it is indicated that he invented and developed a special single-electrode X-ray tube. Tesla's tubes differ from other X-ray tubes in that they have no target electrode. He stated these facts in his 1897 X-ray lecture before the New York Academy of Sciences. The modern term for this process is the bremsstrahlung process, in which a high-energy secondary X-ray emission is produced when charged particles (such as electrons) pass through matter. By 1892, Tesla performed several such experiments; however, he did not categorize the emissions as what was later called X-rays (generalizing the phenonomena as radiant energy). Tesla did not publicly declare his findings nor did he make them widely known. His subsequent X-ray experimentation by vacuum high field emissions led him to alert the scientific community to the biological hazards associated with X-ray exposure. Hermann von Helmholtz formulated mathematical equations for X-rays. He postulated a dispersion theory before Roentgen made his discovery and announcement. It was formed on the basis of the electromagnetic theory of light (Wiedmann's Annalen, Vol. XLVIII); however, he did not work with actual X-rays. On November 8, 1895, Wilhelm Röntgen, a German scientist, began observing and further documenting X-rays while experimenting with vacuum tubes. Röntgen, on December 28, 1895, wrote a preliminary report "On a new kind of ray: A preliminary communication". He submitted it to the Würzburg's Physical-Medical Society journal. This was the first formal and public recognition of the categorization of X-rays. Röntgen referred to the radiation as "X", to indicate that it was an unknown type of radiation. The name stuck, although (over Röntgen's great objections), many of his colleagues suggested calling them Röntgen rays. They are still referred to as Röntgen rays in some countries. Roentgen received the first Nobel Prize in Physics for his discovery. In 1895, Thomas Edison investigated materials' ability to fluoresce when exposed to X-rays. He found that calcium tungstate was the most effective substance. Around March 1896, the fluoroscope he developed became the standard for medical X-ray examinations. Nevertheless, Edison dropped x-ray research around 1903 after the death of Clarence Madison Dally, one of his glassblowers. Dally had a habit of testing X-ray tubes on his hands, and acquired a cancer in them so tenacious that both arms were amputated in a futile attempt to save his life.[1] In 1906, physicist Charles Barkla discovered that X-rays could be scattered by gases, and that each element had a characteristic X-ray. He won the 1917 Nobel prize for this discovery. The use of X-rays for medical purposes was pioneered by Major John Hall-Edwards in Birmingham, England. In 1908, he had to have his left arm amputated owing to the spread of X-ray dermatitis. Detectors The detection of X-rays is based on various methods. Most commonly known is the photographic plate as we know them from hospitals. Initially, most common detection was based on the ionisation of gasses. One of the most simple examples is the Geiger counter. Since the 1990-ies, new detectors were developed based on semiconductors. In the semiconductor, the X-ray photon was converted to electron-hole pairs which were collected to detect the X-ray See also * X-ray crystallography * X-ray Astronomy * X-ray machine * Geiger counter Used to take photographs. Gamma: Above 60 exahertz/Below 5 picometers. Gamma rays (often denoted by the Greek letter gamma, γ) are an energetic form of electromagnetic radiation (see Electromagnetic spectrum) produced by radioactivity or other nuclear or subatomic processes such as electron-positron annihilation. Gamma rays are more penetrating than either alpha or beta radiation, but less ionizing. Gamma rays are distinguished from X rays by their origin. Gamma rays are produced by nuclear transitions while X-rays are produced by energy transitions due to accelerating electrons. Because it is possible for some electron transitions to be of higher energy than nuclear transition, there is an overlap between low energy gamma rays and high energy X-rays. Nuclear processes Radioactive decay processes * Alpha decay * Beta decay * Electron capture * Gamma radiation * Neutron emission * Positron emission * Proton emission * Spontaneous fission Nucleosynthesis * Neutron Capture o The R-process o The S-process * Proton capture: o The P-process Shielding for γ rays requires large amounts of mass. Shields that reduce gamma ray intensity by 50% include 1 cm (0.4 inches) of lead, 6 cm (2.4 inches) of concrete or 9 cm (3.6 inches) of packed dirt. Gamma rays from nuclear fallout would probably cause the largest number of casualties in the event of the use of nuclear weapons in a nuclear war. An effective fallout shelter reduces human exposure at least 1000 times. Gamma rays are less ionising than either alpha or beta rays. However, reducing human danger requires thicker shielding. They produce damage similar to that caused by X-rays such as burns, cancer, and genetic mutations. In terms of ionization, gamma radiation interacts with matter via three main processes: the photoelectric effect, Compton scattering, and pair production. Photoelectric Effect: This describes the case in which a gamma photon interacts with and transfers all of its energy to an orbital electron, ejecting that electron from the atom. The kinetic energy of the resulting photoelectron is equal to the energy of the incident gamma photon minus the binding energy of the electron. The photoelectric effect is thought to be the dominant energy transfer mechanism for x-ray and gamma ray photons with energies below 50 keV (thousand electronvolts), but it is much less important at higher energies. Compton Scattering: This is an interaction in which an incident gamma photon loses enough energy to an orbital electron to cause its ejection, with the remainder of the original photon's energy being emitted as a new, lower energy gamma photon with an emission direction different from that of the incident gamma photon. The probability of Compton scatter decreases with increasing photon energy. Compton scattering is thought to be the principal absorption mechanism for gamma rays in the intermediate energy range 100 keV to 10 MeV (million electronvolts), an energy spectrum which includes most gamma radiation present in a nuclear explosion. Compton scattering is relatively independent of the atomic number of the absorbing material. Pair Production: By interaction in the vicinity of the coulomb force of the nucleus, the energy of the incident photon is spontaneously converted into the mass of an electron-positron pair. A positron is a positively charged electron. Energy in excess of the equivalent rest mass of the two particles (1.02 MeV) appears as the kinetic energy of the pair and the recoil nucleus. The electron of the pair, frequently referred to as the secondary electron, is densely ionizing. The positron has a very short lifetime. It combines within 10−8 second with a free electron. The entire mass of these two particles is then converted into two gamma photons of 0.51 MeV energy each. Gamma rays are often produced alongside other forms of radiation such as alpha or beta. When a nucleus emits an α or β particle, the daughter nucleus is sometimes left in an excited state. It can then jump down to a lower level by emitting a gamma ray in much the same way that an atomic electron can jump to a lower level by emitting ultraviolet radiation. Gamma rays, x-rays, visible light, and UV rays are all forms of electromagnetic radiation. The only difference is the frequency and hence the energy of the photons. Gamma rays are the most energetic. An example of gamma ray production follows. First cobalt-60 decays to excited nickel-60 by beta decay: Then the nickel-60 drops down to the ground state (see nuclear shell model) by emitting a gamma ray: Uses: The powerful nature of gamma-rays have made them useful in the sterilising of medical equipment by killing bacteria. They are also used to kill bacteria in foodstuffs to keep them fresher for longer. In spite of their cancer-causing properties, gamma rays are also used to treat some types of cancer. In the procedure called gamma-knife surgery, multiple concentrated beams of gamma rays are directed on the growth in order to kill the cancerous cells. The beams are aimed from different angles to focus the radiation on the growth while minimising damage to the surrounding tissues. See also: physics, gamma-ray astronomy, gamma ray bursters, radiation therapy. anything above Highly energetic photons associated with black holes and galaxies<br> also emitted by nuclear blasts causes cellular death Black holes are known to emit gamma ray bursts through their poles. A large enough gamma ray burst could sterilize all life on earth.

Electronics/Superconductors. In a superconductor, resistance drops to zero and current flows as a short without attenuating. Until recently most superconductors operated at very low temperatures close to absolute zero. Recently there have been higher temperature superconductors. Superconductors all for floating magnets. Superconductors have a great promise for reducing electricity bills.

Electronics/Motors. Most electrical motors are induction based. They work either through the flow of electricity or turing a magnet. Waterfalls and windmills work by turning a magnet. Waterfalls tend to have even power, but windmills have odd erratic power based on the wind. The driving force that turns the rotor containing the magnets is called the "Prime Mover" Generators work by supplying current that drives a turbine. Induction based motors work on AC and spinning magnets. One of the problems of induction based motors is most of their power is reactive (caught in a magnetic field). Capacitors are required to make the power real. It is common in induction based generators to have a bank of variable capacitors to try to keep the power as real as possible.

Electronics/DC Voltage and Current. Ohm's Law. Ohm's law describes the relationship between voltage, current, and resistance.Voltage and current are proportional to the potential difference and inversely proportional to the resistance of the circuit In this example, the current going through any point in the circuit, "I", will be equal to the voltage "V" divided by the resistance "R". In this example, the voltage across the resistor, "V", will be equal to the supplied current, "I", times the resistance "R". If two of the values ("V", "I", or "R") are known, the other can be calculated using this formula. Any more complicated circuit has an equivalent resistance that will allow us to calculate the current draw from the voltage source. Equivalent resistance is worked out using the fact that all resistors are either in parallel or series. Similarly, if the circuit only has a current source, the equivalent resistance can be used to calculate the voltage dropped across the current source. Kirchoff's Voltage Law. Kirchoff's Voltage Law (KVL): Kirchoff's Current Law. Kirchoff's Current Law (KCL): KCL Example. -"I"1 + "I"2 + "I"3 = 0 ↔ "I"1 = "I"2 + "I"3 "I"1 - "I"2 - "I"3 - "I"4 = 0 ↔ "I"2 + "I"3 + "I"4 = "I"1 Here is more about , which can be integrated here Consequences of KVL and KCL. Voltage Dividers. If two circuit elements are in series, there is a voltage drop across each element, but the current through both must be the same. The voltage at any point in the chain divides according to the resistances. A simple circuit with two (or more) resistors in series with a source is called a voltage divider. Figure A: Voltage Divider circuit. Consider the circuit in Figure A. According to KVL the voltage formula_3 is dropped across resistors formula_4 and formula_5. If a current i flows through the two series resistors then by Ohm's Law. So Therefore Similary if formula_9 is the voltage across formula_4 then In general for n series resistors the voltage dropped across one of them say formula_12 is Where Voltage Dividers as References. Clearly voltage dividers can be used as references. If you have a 9 volt battery and you want 4.5 volts, then connect two equal valued resistors in series and take the reference across the second and ground. There are clearly other concerns though, the first concern is current draw and the effect of the source impedance. Clearly connecting two 100 ohm resistors is a bad idea if the source impedance is, say, 50 ohms. Then the current draw would be 0.036 mA which is quite large if the battery is rated, say, 200 milliampere hours. The loading is more annoying with that source impedance too, the reference voltage with that source impedance is formula_15. So clearly, increasing the order of the resistor to at least 1 kformula_16 is the way to go to reduce the current draw and the effect of loading. The other problem with these voltage divider references is that the reference cannot be loaded if we put a 100 Ω resistor in parallel with a 10 kΩ resistor. When the voltage divider is made of two 10 kΩ resistors, then the resistance of the reference resistor becomes somewhere near 100 Ω. This clearly means a terrible reference. If a 10 MΩ resistor is used for the reference resistor will still be some where around 10 kΩ but still probably less. The effect of tolerances is also a problem; if the resistors are rated 5% then the resistance of 10 kΩ resistors can vary by ±500 Ω. This means more inaccuracy with this sort of reference. Current Dividers. If two elements are in parallel, the voltage across them must be the same, but the current divides according to the resistances. A simple circuit with two (or more) resistors in parallel with a source is called a current divider. Figure B: Parallel Resistors. If a voltage V appears across the resistors in Figure B with only formula_4 and formula_5 for the moment then the current flowing in the circuit, before the division, i is according to Ohms Law. Using the equivalent resistance for a parallel combination of resistors is The current through formula_4 according to Ohms Law is Dividing equation (2) by (1) Similarly In general with n Resistors the current formula_25 is Or possibly more simply Where

Electronics/AC Voltage and Current. Relationship between Voltage and Current. Resistor. In a resistor, the current is in phase with the voltage always. This means that the peaks and valleys of the two waveforms occur at the same times. Resistors can simply be defined as devices that perform the sole function of inhibiting the flow of current through an electrical circuit. Resistors are commercially available having various standard values, nevertheless variable resistors are also made called potentiometers, or pots for short. In theory electricity is a method used to harness myriad numbers including symbols to give the notion on how circuits function on a schematic drawing Capacitor. The capacitor is different from the resistor in several ways. First, it consumes no real power. It does however, supply reactive power to the circuit. In a capacitor, as voltage is increasing the capacitor is charging. Thus a large initial current. As the voltage peaks the capacitor is saturated and the current falls to zero. Following the peak the circuit reverses and the charge leaves the capacitor. The next half of the cycle the circuit runs mirroring the first half. The relationship between voltage and current in a capacitor is: formula_1. This is valid not only in AC but for any function v(t). As a direct consequence we can state that in the real world, the voltage across a capacitor is always a continuous function of the time. If we apply the above formula to a AC voltage (i.e. formula_2), we get for the current a 90° phase shift: formula_3. In an AC circuit, current leads voltage by a quarter phase or 90 degrees. Note that while in DC circuits after the initial charge or discharge no current can flow, in AC circuits a current flows all the time into and out of the capacitor, depending on the in the circuit. This is similar to the in DC circuits, except that the impedance has 2 parts; the resistance included in the circuit, and also the reactance of the capacitor, which depends not only on the size of the capacitor, but also on the of the applied voltage. In a circuit that has DC applied plus a , a capacitor can be used to block the DC, while letting the signal continue. Inductor. In inductors, current is the negative derivative of voltage, meaning that however the voltage changes the current tries to oppose that change. When the voltage is not changing there is no current and no magnetic field. In an AC Circuit, voltage leads current by a quarter phase or 90 degrees. Voltage Defined as the derivative of the flux linkage: formula_4 Resonance. A circuit containing resistors, capacitors, and inductors is said to be in resonance when the reactance of the inductor cancels that of the capacitor to leave the resulting total resistance of the circuit to be equal to the value of the component resistor. The resonance state is achieved by fine tuning the frequency of the circuit to a value where the resulting impedance of the capacitor cancels that of the inductor, resulting in a circuit that appears entirely resistive. See also. AC Voltage: When the polarity of the Voltage between two points is changing continuously with time then the voltage is called AC Voltage. For AC Voltages there will be no constant value so we define it by the average value called RMS value RMS means Root Means Square. In retrospect on AC the power that is used to theoretically determine by an unknown which we assume to be a steady flow, this is also what technicians label the fluctuation of current that determines how electricity moves in a diagram.

Engineering Mechanics. "Engineering mechanics" is the application of mechanics to solve problems involving common engineering elements. The goal of this Engineering Mechanics course is to expose students to problems in mechanics as applied to plausibly real-world scenarios. Problems of particular types are explored in detail in the hopes that students will gain an inductive understanding of the underlying principles at work; students should then be able to recognize problems of this sort in real-world situations and respond accordingly. Further, this text aims to support the learning of Engineering Mechanics with theoretical material, general key techniques, and a sufficient number of solved sample problems to satisfy the first objective as outlined above. Distinction between branches of physics. As you see in the diagram mechanics is the first and most fundamental branch of physics, supporting Thermodynamics and Electricity, and including Statics, Dynamics (= Kinematics + kinetics); all of which are highly applicable in engineering. but the most important part of them is statics (study of body at rest) which is not only a base for all others, but also have the highest engineering application. Physics also involve optics, waves, quantum, and relativity theory, which have no fundamental engineering application yet. Distinction between classic and modern physics. Classical physics generally includes the old theories in physics that are often used to describe the states of physical phenomena in the macroscopic scale e.g. laws governing the laws of motion of masses usually larger than atoms and molecules. Modern physics on the other hand deals with the more recent theories which govern physical phenomena. Two broad branches of modern physics include quantum physics and relativity. These modern theories surfaced as a result of probes into what happened to matter at much more microscopic and abstract states. Hence, classical theories of physics are often classified as theories dating before the 1900s, while modern physics include theories propounded in the 20th and 21st centuries. Prerequisites. This book assumes familiarity with high school physics and calculus, although the mathematics used is fairly elementary. Statics. We describe the motion of bodies using Newton's second law of motion Statics deals with the situation where the acceleration (formula_2) is zero - which happens when a body is at rest, or moving with constant velocity. That means that the total force on a body at rest or with constant velocity must be zero. In other words, the sum of the forces on the body must equal zero. Free body diagram: a tool for understanding and solving. Engineers typically draw what is called a "free body diagram" to show all forces on a body at rest. These forces are then broken down into vectors consistent with a useful coordinate system and summed in sets (components parallel to each basis vector) which are then set to zero to meet the static constraint of no acceleration being present. This typically results in sets of equations which can be solved using simple linear algebra techniques or even simple algebra and substitution. Truss. Forces act along the members, and there are no shear forces or moments. A truss is therefore defined as a system composed entirely of two-force members, which only carry axial loads. The ends of a truss are pinned, so that they don't carry moments. The only reactions at the ends of a truss member are forces. External forces on trusses act only on the end points. Truss problems are solved by the method of sections, where an imaginary cut is made through the member(s) of interest, and global equilibrium of forces and moments are used to determine the forces in the members, or by the method of joints, in which a single joint is isolated and analyzed and the resulting forces (not necessarily with a numerical value) are transferred to adjacent joints, where the process is repeated. The resulting set of equations can then be solved by linear algebra, or substitution. Chains and Cables. Chains and cables are attached at end points and have a continuous load on them due to their own weight (body forces) or external loads. Let a cable of length formula_4 have a load of formula_5 acting per unit distance between the supports. If the tension in the cable at any point is formula_6, then we have, for an infinitesimal length of the cable making an angle formula_7 with the horizontal, Thus, we have, or Now So we have a differential equation for formula_12 in terms of formula_13: Solving for formula_12, If formula_5 is a constant, then we have, which is the equation of a parabola. For a rope where the loading is given in terms of the length of the rope (much more common), i.e., formula_19 and formula_20, we have, where Dynamics. Simply said, "dynamic"s means that an object is in motion with some speed when a force is applied on it. Eg: a car is moving on the road with some velocity. While statics deal with the part of mechanics where all objects are stationary, dynamics deals with objects or structures with a non-zero acceleration. Dynamics may be broken down into kinematics and kinetics. Kinematics deal with displacement, velocities and accelerations without concern for the forces involved. Kinetics deal with the forces and moments involved in making the body move along with the measurement of various parameters describing the motion. Dynamics. Dynamics is a branch of mechanics that deals with the study of bodies in motion. Branches of dynamics. Dynamics is divided into two branches called kinematics and kinetics. Kinematics is a branch of mechanics that describes the movement of a particle or body that does not cause force. Kinetics is a branch of mechanics that combines the force acting on the body with its mass and acceleration.

Abstract Algebra/Rings. This section builds upon and expands the theory covered in the previous chapter on groups. The reader is strongly advised to master the material presented in the sections up to and including Products and Free Groups before continuing. Motivation. The standard motivation for the study of rings is as a generalization of the set of integers formula_1 with addition and multiplication, in order to study integer-like structures in a more general and less restrictive setting. However, we will also present the following motivation for the study of rings, based on the theory of Abelian groups. Let formula_2 and formula_3 be Abelian groups. Then the set formula_4 (Please don't pay much attention to the subscript for now.) of group homomorphisms formula_5 naturally forms an abelian group in the following way. If formula_6, define formula_7 for all formula_8. It should be obvious where each addition is taking place. In particular, we can consider the set formula_9 of endomorphisms of formula_2. That is, the set of homomorphisms from formula_2 to itself. This set is obviously a group from the above discussion, but it is also closed under composition. By endowing the set formula_12 with the operations of addition, formula_13, and composition, formula_14, we note that it has the following properties: Indeed, for the third property, note that if formula_15 and formula_8, then formula_17 and formula_18. The following material is a generalization of this situation. Introduction to Rings. Definition 1: A "ring" formula_19 is a set formula_20 with two binary operations formula_13 and formula_22 that satisfies the following properties: For all formula_23 The definition of ring homomorphism does not include the existence of 1. We will denote the additive identity in a ring by formula_30 or formula_31 if the ring is understood. Similarily, we denote the multiplicative identity by formula_32 or formula_33 when the ring is understood. We'll often use juxtaposition in place of formula_22, i.e., formula_35 for formula_36. Remark 2: Some authors do not require their rings to have a multiplicative identity element. We will call a ring without an idenitity a "rng". "Pseudo-rings" is another term used for rings without unity. Authors who do not require a multiplicative identity usually call a ring a "ring with unity". Unless otherwise stated, we will assume that formula_37 in our rings. A major part of noncommutaive ring theory was developed without assuming every ring has an identity element. Example 3: The reader is already familiar with several examples of rings. For instance formula_38 and formula_39 with the usual addition and multiplication operations. We have a familiy of finite rings given by the sets formula_40 for integer formula_41 with addition and multiplication defined modulo formula_42. Finally we have an example of a "rng" given by the sets formula_43 for integer formula_41 with the usual addition and multiplication. The reader is invited to confirm the ring axioms for these examples. Let us now prove some very basic properties about rings. This is analogous to what we did for groups when we first introduced them. Theorem 4: Let formula_20 be a ring, and let formula_46. Then the following are true: "Proof:" (1), (2), and (3) all strictly concern addition, and are all previous results from formula_24 being a group. The other three parts all concern both addition and multiplication (since 0 and - are additive concepts), so as a proof strategy we expect to use the distributive law in some way to link the two operations. For (4), observe that formula_55. But then by (1), 0a=0. For (5), Note that formula_56. For (6) note that formula_57. Remark 5: Take another look at the examples in Example 3. Notice that for all those rings, multiplication is a commutative opration. However, the axioms say nothing about this. Thus we should expect to find counter-examples to this. Definition 6: A ring is called "commutative" if multiplication is commutative. Example 7: An example of a non-commutative ring is the set formula_58 of formula_59 square matrices with real coefficients under standard addition and multiplication of matrices, where formula_41 is an integer. The reader can easily check this for formula_61 and conclude that it holds for all other formula_42 (why?). Theorem 8: A ring has a unique multiplicative identity. "Proof": During our brief discussion of monoids earlier, we showed that in any monoid the identity is unique. Since a ring sans addition is a monoid, this applies here. Example 9: The singleton set formula_63 with addition and multiplication defined by formula_64 and formula_65 is a ring, called the "trivial ring" or the "zero ring". Note that in the trivial ring, formula_66. The reader is invited to show that formula_66 in a ring if and only if it is the trivial ring. If the reader has tried to construct some of the rings formula_68, he/she may have realised that certain non-zero elements have product zero. We formalize this concept as follows. Definition 10: Let formula_20 be a ring and formula_70. formula_71 is called a "left(resp.right)-zero-divisor" if there exists a formula_72 such that formula_73 formula_74. Lemma 11: Let formula_20 be a ring with formula_76. Define the function formula_77 given by formula_78 for all formula_79. Then formula_80 is injective if and only if formula_71 is not a left-zero-divisor. "Proof": Assume formula_71 is not a left-zero-divisor, and assume we have formula_83 for some formula_84. This implies formula_85, giving formula_86 since formula_71 is not a left-zero-divisor, so formula_80 is injective. Conversely, assume formula_71 is a left-zero-divisor. Then there exists a formula_72 such that formula_91 and formula_92, so formula_80 is not injective. Remark 12: Thus, multiplication by formula_71 is left-cancellative if and only if formula_71 is not a zero-divisor. The reader is invited to state and prove the equivalent lemma for right-zero-divisors. Example 13: formula_96 are all examples of commutative rings without zero divisors. These rings motivate the next definition. Definition 14: Let formula_20 be a commutative ring without zero divisors. Then formula_20 is called an "integral domain". Just like Definition 14, the majority of special types of rings will be motivated by properties of formula_1. Example 15: Definition 16: Let formula_20 be a ring. An element formula_111 is a "unit" and is "invertible" if there is an element formula_112 such that formula_113. The set of all units is denoted by formula_114. Exercise 17: Prove that formula_114 is a group under multiplication. Exercise 18:: Show that a zero-divisor is not a unit. Theorem 19: (Cancellation Law for Integral Domains): Let formula_20 be an integral domain, and let formula_46 be nonzero. Then formula_118 if and only if formula_48. "Proof:" Evidently formula_118 if formula_86. To see the other direction, we rearrange the equality as formula_122. But then formula_123. Since formula_71 is nonzero, and formula_20 contains no zero divisors, it must be the case that formula_126, which is to say that formula_86. Definition 20: A ring formula_20 is a "division ring" or "skew field" if all non-zero elements are units, i.e. if it forms a group under multiplication with its nonzero elements. Definition 21: A field is a commutative division ring. Alternatively, a field formula_129 is a ring where formula_130 is an abelian group under multiplication. As another alternative, a field is an integral domain where all non-zero elements are invertible. As stated before, integral domains are easy to work with because they are so close to being fields. In fact, the next theorem shows just how close the two are: Theorem 22: Let formula_20 be a finite integral domain. Then formula_20 is a field. "Proof:" Let formula_76 be nonzero and let formula_134. Clearly formula_135 is a subset of formula_20. From the cancellation law, we can see that formula_137 (since if two elements formula_138 and formula_139 are equal, then formula_86). But then formula_141. So then there must be some formula_142 such that formula_143. So formula_71 is a unit. Of course proving that a set with two operations satisfy all of the ring axioms can be tedious. So, just as we did for groups, we note that if we're considering a subset of something that's already a ring, then our job is easier. Definition 23: A "subring" formula_135 of a ring formula_20 is a subset of formula_20 that is also a ring (under the same two operations as for formula_20) and formula_149. We denote "formula_135 is a subring of formula_20" by formula_152. Note many mathematicians do not require rings or subrings to have an identity. Theorem 24: Let formula_153 be a subset of a ring formula_20. Then formula_152 if and only if for all formula_156, Example 25:

German/Level I/Introduction. I.0: Introduction Welcome to Level I German! Level I is aimed at junior high and high school students. However, it can also be used by others just beginning to learn to speak or read German. The goal of Level I German is to introduce the basics of the German language without overwhelming students. Therefore, the vocabulary is formatted for translating from English (which the students know) into German. Although Level II is aimed at students and people who are a bit proficient after Level I, still, English translation will be used, so as to ease the learning. It helps because, at times while learning a new language, even with basic understanding, the words are above normal understanding level, and thus require a "sub" assistance. German and English. German and English are quite close to each other, and are called language sisters or, more formally, "cognate languages". Both belong to the Germanic branch of the Indo-European language family. Here are some major similarities: As you can see, German is quite similar to English. There are, however, differences: However, German is still one of the easiest languages for English speakers to learn. The differences will be tackled over the course of the lessons. How to Use this Level of the German Textbook. The lessons are meant to be taken in order. You should read and review the German dialogs as often as possible. Many of the dialogs come with audio recordings by native speakers. These recordings are invaluable to learn the German pronunciation. If there is a recording, you can do several kinds of exercises: At the reviews, after every third lesson, go back to look at the previous lessons. Levels of Completion. On the contents page, you will see filled-in boxes next to each lesson. The number of boxes corresponds to the completeness of the lesson as follows:

Trigonometry/For Enthusiasts/Less-Used Trig Identities. Triangle Identities. In addition to the Law of Sines, the Law of Cosines, and the Law of Tangents, there are numerous other identities that apply to the three angles A, B, and C of any triangle (where A+B+C=180° and each of A, B, and C is greater than zero). Some of the most notable ones follow: Pythagoras. These are all direct consequences of Pythagoras's theorem. Multiple angle. These are all direct consequences of the sum/difference formulae Half angle. In cases with formula_49 , the sign of the result must be determined from the value of formula_50 . These derive from the formula_51 formulae.

Electronics/Batteries. Electrolysis involves an electric field and the motion of electrons through a liquid. GCSE Science/Electricity Electrolysis is the decomposition of certain types of substance using electricity. The types of substance that can be split are ionic substances. This just means that they are made of charged ions rather than neutral atoms. {Remember that an ion is just an atom that has either a positive or negative charge}. An example of an ionic substance is common table salt—sodium chloride. The sodium ion has a positive charge, the chlorine ion has a negative charge. It is usually written as Na+Cl-. Q1) Check in a periodic table, what is the symbol for sodium—Na or Cl? As you may already know if you've studied the Metals module, a salt is any substance made by combining an acid with an alkali. Acids, alkalis, and therefore all salts are ionic. Q2) Which of the following substances can be broken up by electricity Sodium chloride, iron sulphate, copper nitrate? Basic Experimental Setup. Most ionic compounds are not liquid at room temperature. This is a problem because the ions need to be able to move for the electric current to be able to flow. This can be achieved by melting. Look at the electrical setup shown on the right. The electrodes are just two carbon rods connected to a battery. The one connected to the positive electrode is called the anode. The one connected to the negative electrode is called the cathode. Consider for example the compound lead bromide. This compound is a solid at room temperature but can be molten over a Bunsen flame. So what you would do is put some lead bromide into a beaker. Put the beaker on a tripod over a Bunsen flame. Melt the lead bromide, then put in the electrodes and turn the power supply on at a setting of say 2V. What you would see happening is the anode, is a silvery coating of pure lead forming, and bromine forming at the cathode. The current would continue to flow until all the lead bromide was turned into lead and bromine. Q3) It takes energy to split up a compound like lead bromide. Where does this energy come from ? Q4) Predict what products you would get at the anode and cathode if copper chloride was the electrolyte. What happens at the anode. The anode is the positive electrode. It attracts negatively charged ions, because unlike charges attract. The bromine ions move through the melt until they reach the anode. Once they get there, they give up their two extra electrons to become bromine atoms. Br2- --> Br + 2e- The electrons flow up the anode to the positive terminal of the battery. What happens at the cathode. The cathode is the negative electrode. It attracts the positively charged ions. Metal ions are always positive and so the lead ions flow through the melt until they get to the cathode. Once they get there they take two electrons each from the cathode. The electrons flow down the cathode from the negatively terminal of the battery and onto the lead ions. Pb2+ + 2e- --> Pb Q5) Solid ionic substances do not conduct electricity and are not split up by it. Why do you think that is? Quantity calculations (higher tier only). In the experiment with lead bromide, you saw that lead was deposited at the cathode. If you actually do the experiment you will see that the lead coats the cathode. In this section we will look at how much metal will coat a cathode in a given time. A scientist performed the following experiment. His results were: You can see from the results that the total amount of copper deposited depends on both the current and the time it flows. This is because the number of copper atoms that can be made from ions depends on the total amount of charge that flows. The unit of charge is the coulomb. One coulomb is the amount of charge when one Amp flows for one second (see GCSE Science/Measuring electricity Q6)Look at the results table above. How much copper is deposited when 1A flows for 3000 seconds? Q7)How much copper do you predict would be doposited if 1A were to flow for 6000 seconds. Q8)What about if 2A were to flow for 12000 seconds ? Electrolysis of Aqueos solutions (Advanced). "Before studying this section check with your teacher to see if you need to". Earlier on in this module you've learned that ions must be able to move in order for electrolysis to work. If the ions are held rigid {such as in a solid}, they can't move and no electricity will flow. We've looked at how the freeing up of ions can occur by melting the electrolyte. Another way to achieve this is by dissolving th electrolyte in water. The trouble with this method is, there will be mor than one type of ion present. water partially splits up into ions {this is why it's such a good solvent for ionic compounds}. It splits into hydrogen ions and hydroxide ions. H2O --> H+ +OH- So at the cathode there will be two ions present. The Metal ion and the hydrogen ion from the water. Which element is actually produced at the cathode depends on how reactive the metal is. If the metal is very reactive, such as potassium, or sodium, then hydrogen will be produced. If the metal is unreactive such as silver, the metal will be produced.{I hope you learned your displacement reactions in lower school}. To work out which ion wins, the metal or the hydrogen, compare their reactivities in a the reactivity series. The one that is most reactive, will not be produced at the cathode. A similar situation occurs at the anode. Hydroxide ions {from the water} are usually discharged at the anode ultimately producing oxygen. OH- --> OH + e- 4OH --> 2H2O + O2. Q9) Sodium chloride is dissolved in water and subjected to electrolysis. Explain what you see at each of the electrodes. Answers | «Advanced static electricity | Circuits» A fuel cell is an electrochemical device similar to a battery, but differing from the latter in that it is designed for continuous replenishment of the reactants consumed; i.e. it produces electricity from an external fuel supply as opposed to the limited internal energy storage capacity of a battery. Typical reactants are hydrogen on the anode side and oxygen on the cathode side. In contrast, conventional batteries consume solid reactants and, once these reactants are depleted, must be discarded, recharged with electricity by running the chemical reaction backwards, or, at least in theory, having their electrodes replaced. Typically in fuel cells, reactants flow in and reaction products flow out, and continuous long-term operation is feasible virtually as long as these flows are maintained. Fuel cells are also attractive in some applications for their high efficiency and low pollution. Some applications that have been suggested include Types of fuel cells. There are five generally recognised types of fuel cells, of which two are the main subject of intensive research. PEM. PEM fuel cells have a disputed acronym, meaning either proton-exchange membrane or polymer-electrolyte membrane, which both are in fact a good description. In this fuel cell, hydrogen is split at the membrane surface in protons, that travel through the membrane, and electrons, that travel through our external electric circuit, and provide our power. The hydrogen ions travel through the water that is entrained in the membrane to the other side, where they are combined with oxygen to form water. Unfortunately, while the splitting of the hydrogen molecule is relatively easy, splitting the stronger oxygen molecule is more difficult, and this causes significant losses that result in a sharp decrease of performance of the fuel cell. The PEM fuel cell is a prime candidate for vehicle and other mobile applications of all sizes down to mobile phones, because of its compactness. However, the water-entraining membrane is crucial to performance: too much water will flood the membrane, too little will dry it; in both cases, power output will drop; water management is a very difficult subject in PEM fuel cells. Furthermore, the platinum catalyst on the membrane is easily poisoned by CO. PEM systems that use reformed methanol were proposed, as in Daimler Chrysler Necar 5; reforming methanol, i.e. making it react to obtain hydrogen, is however a very complicated process, that requires also purification from the CO the reaction produces. A platinum-ruthenium catalyst is necessary as some CO will unavoidably reach the membrane. The level should not exceed 10 parts per million. DMFC. A subcategory of PEM is the DMFC, or direct methanol fuel cell; here, the methanol is not reformed, but fed directly to the fuel cell. One does not need complicated reforming, and storage of methanol is much easier than that of hydrogen. However, efficiency is low, due to the high permeation of methanol through the membrane, and the dynamic behaviour is sluggish. The main manufacturer of PEM is Ballard Power Systems of Vancouver, Canada. Efficiencies of PEM are in the range of 40-50%. SOFC. Solid oxide fuel cells, or SOFC, are intended mainly for stationary applications (power plants). They work at very high temperatures (some at 1000°C), and their off-gases can be used to fire a secondary gas turbine to improve efficiency. Efficiency could reach as much as 70% in these hybrid systems. This time, it's oxygen being transferred through a solid oxide at high temperature to react with hydrogen on the other side. SOFC have such a high temperature that they can be fed (provided some modifications) with natural gas, that will react to give hydrogen in the fuel cell itself. SOFC are very resistant to poisoning, and can indeed be run on CO, which is a poison for PEM. Since SOFC are made of ceramic materials, they tend to be brittle; they are therefore unsuited for mobile applications. Furthermore, thermal expansion demands a uniform and slow heating process at startup, that will cause very long startup times: typically, 8 hours are to be expected. Research is going now in the direction of lower-temperature SOFC (600°C), which will enable the use of metallic materials with better mechanical properties and heat conductivity. MCFC. Molten-carbonate fuel cells (MCFC) are also high-temperature, but in the range of 600°C. Their main problem is corrosion, and the need to operate a high-temperature liquid rather than a solid as in the SOFC. PAFC. Phosphoric-acid fuel cells (PAFC) are a mature technology that is commercially available. Unfortunately, the phosphoric acid solidifies under 40°C, making startup very difficult. They have been used for stationary applications with an efficiency of about 40%, and many believe they do not offer much potential for further development. AFC. The alkaline fuel cell (AFC) is the cell that brought the Man to the Moon. Used in Apollo-series missions and on the Space Shuttle, it is a very good fuel cell but for the fact that it is poisoned by CO2. This means that the cell will require pure oxygen, or at least purified air. As this process is relatively expensive, not much development is being done on AFC. NASA has decided they will shift to PEM for the next generation of Space Shuttles. Science. Fuel cells are electrochemical devices, so they are not constrained by the maximum thermal (Carnot) efficiency as combustion engines are. Consequently, they can have very high efficiencies in converting chemical energy to electrical energy. In the archetypal example of a hydrogen/oxygen polymer electrolyte membrane (PEM) fuel cell, a proton-conducting polymer membrane separates the anode ("fuel") and cathode sides. Each side has an electrode, typically carbon paper coated with platinum catalyst. On the anode side, hydrogen diffuses to the anode catalyst where it dissociates into protons and electrons. The protons are conducted through the membrane to the cathode, but the electrons are forced to travel in an external circuit (supplying power) because the membrane is electronically insulating. On the cathode catalyst, oxygen molecules react with the electrons (which have travelled through the external circuit) and protons to form water. In this example, the only waste product is water vapor. Also, there is the possible use of fuel cells at home, to store energy at the cheap off-peak electricity rates and used at peak-use hours. It may even be profitable to sell back some of the energy to the power company, like they do with windmill electric power. Peak power production reaches twice the normal level, which means that the very expensive powerplant capacity is sized for levels used for a short period of time. Also, power plants are most efficient at only one production rate and their efficiency drops off significantly at off-peak rates. History. The first fuel cell was developed in the 19th century by British scientist Sir William Grove. A sketch was published in 1843. But fuel cells did not see practical application until the 1960s, where they were used in the U.S. space program to supply electricity and drinking water (hydrogen and oxygen being readily available from the spacecraft tanks). Extremely expensive materials were used and the fuel cells required very pure hydrogen and oxygen. Early fuel cells tended to require inconveniently high operating temperatures that were a problem in many applications. Further technological advances in the 1980s and 1990s, like the use of Nafion as the electrolyte, and reductions in the quantity of expensive platinum catalyst required, have made the prospect of fuel cells in consumer applications such as automobiles more realistic. The fuel cell industry. Ballard Power Systems is a major manufacturer of fuel cells and leads the world in automotive fuel cell technology. Ford Motor Company and Chrysler are major investors in Ballard. As of 2003, the only major automobile companies pursuing internal development of fuel cells for automotive use are General Motors and Toyota; most others are customers of Ballard. United Technologies (UTX) is a major manufacturer of large, stationary fuel cells used as co-generation power plants in hospitals and large office buildings. The company has also developed bus fleets that are powered by fuel cells. Pros and cons of fuel cells in various applications. Their use is controversial in some applications. The hydrogen typically used as a fuel isn't a primary source of energy. It is usually only a source of stored energy that must be manufactured using energy from other sources. Some critics of the current stages of this technology argue that the energy needed to create the fuel in the first place may reduce the ultimate energy efficiency of the system to below that of highly efficient gasoline internal-combustion engines; this is especially true if the hydrogen is generated from electrolysis of water by electricity. On the other hand, hydrogen can be generated from methane (the primary component of natural gas) with approximately 80% efficiency. The methane conversion method releases greenhouse gases, however, and the ideal environmental system would be to use renewable energy sources to generate hydrogen through electrolysis. Other types of fuel cells don't face this problem. For example, biological fuel cells take glucose and methanol from food scraps and convert it into hydrogen and food for the bacteria. There are practical problems to be overcome as well. Although the use of fuel cells for consumer products is probable in the near future, most current designs won't work if oriented upside down. They currently cannot be scaled to the small size needed by portable devices such as cell phones. Current designs require venting and therefore cannot operate under water. They may not be usable on aircraft because of the risk of fuel leaks through the vents. Technologies for safe refueling of consumer fuel cells are not yet in place. Among the controversies in the use of hydrogen are: First, the energy used to produce the hydrogen is comparable to the energy in the hydrogen, it is inefficient therefore too expensive. If conventional powerplants were used to produce the hydrogen, at best, there would be no gains in current pollution rates. Second, some have suggested that this is a "stalking horse" to bring back nuclear power, which may be the only commercially viable way to make hydrogen and reduce pollution. Finally, the need to provide the very long, costly and vulnerable thousands of miles of gas lines makes hydrogen too costly without government help. There are several advantages to hydrogen as well. Clean, renewable energy sources like solar and wind power are non-continuous and unreliable through the course of a day. So power from these sources is not always available at the time it is needed. The electricity produced from solar panels or wind generators could be stored in large battery complexes but this can be expensive and batteries have a limited storage capacity and lifetime. If the electricity is used to produce hydrogen however, the energy can be stored more easily. As a gas hydrogen is not hard to store until it is needed. Batteries in circuits. Batteries in series can increase or decrease the combined voltage. Toys use batteries, for example, and it is important to note from the associated drawing(s) which way the batteries are to be inserted into the toy because their voltages can add or subtract as the case may be. To add voltages the + of one battery is connected to the - of the next battery. Some damage may be caused by inserting batteries in the wrong direction. Batteries in parallel are not usually used; at a given rated voltage each battery is able to provide only a certain maximum current, but to increase the maximum available output current several batteries could be connected in parallel, meaning that all the + terminals are connected together, and so are all of the - terminals. The life of a battery depends approximately on the value of (current times time), but the chemistry involved can permit a battery that seems to have been finished its lifetime to regain some of its life after some time, after the chemistry was able to do its work. Many different kinds of batteries exist. They are used for different purposes, using different materials, sizes, rated voltages. Types of batteries. Reversible, rechargeable, irreversible. Alkaline = Basic. AA<br> AAA<br> C<br> D<br> 9V<br> All batteries, except the 9V variety, are 1.5V. As alkaline batteries lose their charge, or in Layman's terms "juice", the voltage outputted slowly decreases. This is what causes flashlights to dim and remote controls to lose their effectiveness whilst using the last of the power left in the battery. Rechargeable. Lithium. Lithium batteries are available in rechargeable and "normal" forms. Rechargeable lithium-ion batteries are typically found in small to medium electronics, such as laptops, music players, some cordless mice etc. Lithium batteries have little/no proven research of memory effect. Lithium batteries are also available as single decharge cycle batteries. These batteries are marketed as using lithium to increase life span but offer little/no fiscal advantage over standard alkaline batteries. Nickel Cadmium (NiCd). Nickel Cadmium batteries are typically not used in the present but were one of the first easily obtainable cheap rechargeable batteries besides the battery. Nickel Cadmium (AKA NiCd) batteries can typically be found in radio controlled cars, older rechargeable "standard" batteries (AA,AAA, etc.), some wireless mice and many other applications. There is a myth that Nickel-Cadmium batteries have a memory effect, where if they are recharged when still somewhat full, they will lose capacity. laboratory research has proven this false in modern batteries (it does happen, to a small extent in old (pre-1950) sintered-plate NiCds). NiCds do suffer from voltage depression, which is similar to the memory effect, but only occurs when the batteries are overcharged. Nickel Metal Hydride (NiMH). NiMH batteries can be looked at as a newer version of NiCd batteries. NiMH batteries are often used in the same applications as NiCd batteries. However, they last longer and don't have voltage depression as severe as NiCd batteries. Lead Acid. Car and boat batteries are the most common forms of lead acid batteries. In both uses they are typically rated for 12V. They have a high current output and can severely injure a person who was to short the circuit. Memory effect. The memory effect is a myth that states when the battery is not fully charged before it is recharged, it loses capacity, it "forgets" that it has the capacity. The myth also states that the memory effect can sometimes be reversed by using a conditioner. A battery conditioner almost completely discharges a battery (down to 1 volt in NiCds) and are designed to reverse voltage depression, not memory effect. the memory effect does not exist except in very old (1950 or earlier) sintered plate NiCds. There is no need to use a battery conditioner on any battery to reverse memory effect.

Indonesian/Lesson 0 Exercise. < Indonesian/Lesson 0 Okay, so now you at least know how to read Indonesian. That's a great start! Indonesian/Malay has a rhymic quality akin to Italian. And for long, native poets, scribes, and orators have had a penchant of creating sentences that rhyme with each other. Just like in English poetry, the rhyming makes extensive use of the many similar final-sounds of words. To prove my point, let me bring you the Rhyme Game. Rhyme Game. Objective: Kuku-kukuku kaku (My nails are stiff). -ku (pronounced -COO) is not very common, but it does contain some important words! Now that I've established the pattern, can you guess how to pronounce these words? Cinta itu buta (Love is blind). Words ending with -ta (-TUH) are often related to love, and to umm... many other things! (*)the vowel sound in ren- is essentially a schwa, in contrast to gEn- in genta Can you guess, based on what you've learned? Tangkap ikan kakap! (Catch the Barramundi fish!). When you catch and hold onto something, you say -kap (-CUP)! And when you let something be known, you should also say -kap! Here we go again! Other Tongue Twisters. If you're feeling especially game after your Indonesian phonetics lesson, have a go at the following tongue twisters:

Latin/Introduction to Nouns (L1). =Introduction to Nouns= A noun in Latin is exactly like a noun in English - it is a person, place, thing, or idea. For more information on what a noun is and what words are nouns look at "" at Wikipedia. If something is a noun in English, then it is also a noun in Latin. =Nouns in Latin= Okay, now that we know that nouns are essentially the same part of speech in both English and Latin, what makes Latin nouns work? Anytime that you see a noun in Latin, you will be able to tell a lot about the noun. For example, if you were to see the noun "mensam" in a Latin text, you would know a number of things about it immediately. We'll go into how you can tell these things in just a bit, but here's what you can learn from seeing the word. First, the word is almost certainly feminine. Second, it's probably the direct object of a sentence. Third, its singular. To compare that to English, let's look at the English word "table" (because that's what "mensam" means!) We can see that it is singular, because the plural is "tables." But what else can we learn about the word? It isn't necessarily a direct object--it could be a subject (The table is large), indirect object (I walked to the table), or direct object (He kicked the table). Also, we don't know its gender because English doesn't use a gender system for nouns. So, why are we able to see these things about Latin words? We know them because Latin nouns "decline" - they change their endings to fit their usage. "Mensam" and "mensa" both mean table, but they are used very differently. The difference in the form of a noun gives us information about the word. When we talk about the form a word in Latin takes, we talk about its "incidents". English nouns typically only have one incident - number (singular vs plural). However, Latin nouns have three: Gender. All nouns in Latin have gender. This includes not only things which physically have gender (man, woman, boy, girl, mother, husband, etc.) but also inanimate objects like tables. "Mensa" (table), for example, is feminine just like "puella" (girl). Sometimes a thing (usually animals or words for people like "student" will have two different nouns, with one masculine and the other feminine. "Lupus" is a male wolf while "lupa" is a female wolf. "Discipula" is a female student while "discipulus" is a male student. What really throws a wrench into this is that Latin does not have two genders. It has three. In addition to masculine and feminine, Latin also has the neuter gender (interestingly enough, this comes from the Latin "neuter" which means "neither"). "Crustulum", a cookie, is a good example of a neuter noun. One might ask, why is a table feminine while a cookie is neuter? The simplest answer is: because that's just the way it is. Sometimes the gender of Latin nouns will make absolutely no sense. At other times, it will. It is best, however, to memorize the gender of each noun you learn. Case. The next step in analyzing a Latin noun is its case. The case of a noun determines its usage in a sentence. Each noun can be put into any of the cases. If you look above, you can see that you have already seen two different forms of the Latin word for "table". They are in different cases. English for the most part does not use cases, but some pronouns show vestiges of a case system ("who" vs. "whom" vs. "whose" and "they" vs. "them" vs. "their"). We take a fairly intuitive approach to learning those words in English, which is perhaps why so many people have problems with "whom". But in Latin, the case of a noun is extremely important. The number of cases in Latin is generally considered to be six. This leaves out the locative, which is used very rarely and only for a very limited number of nouns, so we'll leave it out for now as well. We'll start off with the standard six - nominative, accusative, genitive, dative, ablative, and vocative. Note that all forms in the above chart are singular (we'll talk about number in just a bit). You can see that sometimes a couple different cases will look the same. This happens. It is usually possible to determine the appropriate case from the context of a sentence (for instance, you could say "I walked to the table" but wouldn't say "I walked of the table"). Each noun belongs to a group called a "declension" - all of the nouns in a declension have similar sets of endings. For example, "mensa" is a first declension noun. All of the other first declension nouns are going to have similar endings. Let's take a look at "mensa" alongside "puella". We'll explore each of the cases and each of the declensions later. Number. This should be fairly familiar to everyone (well, except perhaps speakers of Japanese and some other languages, whose nouns do not distinguish between singular and plural). Number refers only to whether a noun is singular (there's only one of it) or plural (there's more than one). Let's take a look at our table again. In English, we can easily see the difference between "table" and "tables" by the ending. Latin is the same. In the accusative case, "mensam" is singular ("table") while "mensas" is plural (tables). The ending is determined by the number. However, this can create even more confusion. For example, you can see by the chart above (which shows singular forms) that "mensae" can be either the genitive singular or dative singular. But guess what? "Mensae" can also be the nominative or vocative plural! Obviously, we need a new chart that covers number as well as case. You can see that "mensae" appears four times, "mensa" appears three times, and "mensis" appears twice. Again, context is necessary in determining which case and number a noun will be in. Dictionary form of a noun. The dictionary form of a noun includes three things. So, if we wanted to look up "table" in our handy Latin-to-English dictionary, we would look under: mensa, mensae, "f". =Summary= Now you have a working knowledge of what constitutes a Latin noun. You know that each noun belongs to one the five declensions and one of the three genders. You know what the six cases are. You're ready to go onto the next chapter and learn about the first declension and the nominative and accusative cases. Latin | Lesson 2 - First Declension, Nominative/Accusative Cases

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XML - Managing Data Exchange/XHTML. In previous chapters, we have learned how to generate HTML documents from XML documents and XSL stylesheets. In this chapter, we will learn how to convert those HTML documents into valid XHTML. We will discuss why XHTML has evolved as a standard and when it should be used. The Evolution of XHTML. Originally, Web pages were designed in HTML. Unfortunately most implementations of this markup language allow all sorts of mistakes and bad formatting. Major browsers were designed to be forgiving, and poor code would display with few problems in most cases. This poor code was often not portable between browsers, e.g. a page would render in Netscape but not Internet Explorer or vice versa. The accounting for human error and bad formatting takes an amount of processing power that small handheld devices might not have. Thus when displaying data on handhelds, a tiny mistake can crash the device. XHTML partially mitigates these problems. The processing burden is reduced by requiring XHTML documents to conform to the much stricter rules defined in XML. Aside from the stricter rules, HTML 4.01 and XHTML 1.0 are functionally equivalent. If a document breaks XML's well-formedness rules, an XHTML-compliant browser must not render the page. If a document is well-formed but invalid, an XHTML-compliant browser may render the page, so a significant number of mistakes still slip through. In this chapter, we will examine in detail how to create an XHTML document. The biggest problem with HTML from a design standpoint is that it was never meant to be a graphical design language. The original version of HTML was intended to structure human readable content (e.g. marking a section of text as a paragraph), not to format it (e.g. this paragraph should be displayed in 14pt Arial). HTML has evolved far past its original purpose and is being stretched and manipulated to cover cases that the original HTML designers never imagined. The recommended solution is to use a separate language to describe the presentation of a group of documents. Cascading Style Sheets (CSS) is a language used for describing presentation. From version 1.1 of XHTML upwards web pages must be formatted using CSS or a language with equivalent capabilites such as XSLT (XSL Transformations). The use of CSS or XSLT is optional in XHTML 1.0 unless the strict variant is used. HTML 4.01 supports CSS but not XSLT. So What is XHTML? As you might have guessed, XHTML stands for eXtensible HyperText Markup Language. It is a cross between HTML and XML. It fulfills two major purposes that were ignored by HTML: The best thing about XHTML is that it is almost the same as HTML! If you know how to write an HTML document, it will be very simple for you to create an XHTML document without too much trouble. The biggest thing that you must keep in mind is that unlike with HTML, where simple errors like missing a closing tag are ignored by the browser, XHTML code must be written according to an exact specification. We will see later that adhering to these strict specifications actually allows XHTML to be more flexible than HTML. XHTML Document Structure. At a minimum, an XHTML document must contain a DOCTYPE declaration and four elements: html, head, title, and body: <!DOCTYPE ... > <html xmlns="http://www.w3.org/1999/xhtml" xml:lang="..."> <head> <title></title> </head> <body></body> </html> The opening codice_4 tag of an XHTML document must include a namespace declaration for the XHTML namespace. The DOCTYPE declaration should appear immediately before the html tag in an XHTML document. It can follow one of three formats. XHTML 1.0 Strict. <!DOCTYPE html PUBLIC "-//W3C//DTD XHTML 1.0 Strict//EN" "http://www.w3.org/TR/xhtml1/DTD/xhtml1-strict.dtd"> The Strict declaration is the least forgiving. This is the preferred DOCTYPE for new documents. Strict documents tend to be streamlined and clean. All formatting will appear in Cascading Style Sheets rather than the document itself. Elements that should be included in the Cascading Style Sheet and not the document itself include, but are not limited to: There are also certain instances where your code needs to be nested within block elements. Incorrect Example: Correct Example: XHTML 1.0 Transitional. This declaration is intended as a halfway house for migrating legacy HTML documents to XHTML 1.0 Strict. The W3C encourages authors to use the Strict DOCTYPE for new documents. (The XHTML 1.0 Transitional DTD refers readers to the relevant note in the HTML4.01 Transitional DTD.) This DOCTYPE does not require CSS for formatting; although, it is recommended. It generally tolerates inline elements found where block-level elements are expected. There are a couple of reasons why you might choose this DOCTYPE for new documents. XHTML 1.0 Frameset. If you are creating a page with frames, this declaration is appropriate. However, since frames are generally discouraged when designing Web pages, this declaration should be used rarely. XML Prolog. Additionally, XHTML authors are encouraged by the W3C to include the following processing instruction as the first line of each document: <?xml version="1.0" encoding="UTF-8"?> Although it is recommended by the standard, this processing instruction may cause errors in older Web browsers including Internet Explorer version 6. It is up to the individual author to decide whether to include the prolog. Language. It is good practice to include the optional codice_5 attribute on the html element to describe the document's primary language. For compatibility with HTML the codice_6 attribute should also be specified with the same value. For an English language document use: <html xmlns="http://www.w3.org/1999/xhtml" xml:lang="en" lang="en"> The codice_5 and codice_6 attributes can also be specified on other elements to indicate changes of language within the document, e.g. a French quotation in an English document. Converting HTML to XHTML. In this section, we will discover how to transform an HTML document into an XHTML document. We will examine each of the following rules: Documents must be well-formed. Because XHTML conforms to all XML standards, an XHTML document must be well-formed according to the W3C's recommendations for an XML document. Several of the rules here reemphasize this point. We will consider both incorrect and correct examples. Tags must be properly nested. Browsers widely tolerate badly nested tags in HTML documents. <b><u> This text is probably bold and underlined, but inside incorrectly nested tags. </b></u> The text above would display as bold and underlined, even though the end tags are not in the proper order. An XHTML page will not display if the tags are improperly nested, because it would not be considered a valid XML document. The problem can be easily fixed. <b><u> This text is bold and underlined and inside properly nested tags. </u></b> Elements must be closed. Again, XHTML documents must be considered valid XML documents. For this reason, all tags must be closed. HTML specifications listed some tags as having "optional" end tags, such as the and <li> tags. <p>Here is a list: <ul> <li>Item 1 <li>Item 2 <li>Item 3 </ul> In XHTML, the end tags must be included. <p>Here is a list: </p> <ul> <li>Item 1</li> <li>Item 2</li> <li>Item 3</li> </ul> What should we do about HTML tags that do not have a closing tag? Some special tags do not require or imply a closing tag. <img src="titlebar.gif" alt="Title"> <hr> <br> <p>Welcome to my web page!</p> In XHTML, the XML rule of including a closing slash within the tag must be followed. <img src="titlebar.gif" alt="title" /> <hr /> <br /> <p>Welcome to my Web page!</p> Note that some of today's browsers will incorrectly render a page if the closing slash does not have a space before it (). Although it is not part of the official recommendation, you should always include the space () for compatibility purposes. Here are the common empty tags in HTML: Tags must be lowercase. In HTML, tags could be written in either lowercase or uppercase. In fact, some Web authors preferred to write tags in uppercase to make them easier to read. XHTML requires that all tags be lowercase. <H1>This is an example of bad case.</h1> This difference is necessary because XML differentiates between cases. XML would read and as different tags, causing problems in the above example. <h1>This is an example of good case.</h1> The problem can be easily fixed by changing all tags to lowercase. Attribute names must be lowercase. Following the pattern of writing all tags in lowercase, all attribute names must also be in lowercase. <p CLASS="specialText">Important Notice</p> The correct tags are easy to create. <p class="specialText">Important Notice</p> Attribute values must be quoted. Some HTML values do not require quotation marks around them. They are understood by browsers. <table border=1 width=100%> </table> XHTML requires all attributes to be quoted. Even numeric, percentage, and hexadecimal values must appear in quotations for them to be considered part of a proper XHTML document. <table border="1" width="100%"> </table> Attributes cannot be minimized. HTML allowed some attributes to be written in shorthand, such as selected or noresize. <form> <input checked ... /> <input disabled ... /> </form> When using XHTML, attribute minimization is forbidden. Instead, use the syntax x="x", where x is the attribute that was formerly minimized. <form> <input checked="checked" .../> <input disabled="disabled" .../> </form> A complete list of minimized attributes follows: The codice_9 attribute is replaced with the codice_10 attribute. HTML 4.01 standards define a name attribute for the tags a, applet, frame, iframe, img, and map. <a name="anchor"> <img src="banner.gif" name="mybanner" /> </a> XHTML has deprecated the name attribute. Instead, the id attribute is used. However, to ensure backwards compatibility with today's browsers, it is best to use both the name and id attributes. <a name="anchor" id="anchor" > <img src="banner.gif" name="mybanner" id="mybanner" /> </a> As technology advances, it will eventually be unnecessary to use both attributes and XHTML 1.1 removed name altogether. Ampersands are not supported. Ampersands are illegal in XHTML. <a href="home.aspx?status=done&itWorked=false">Home & Garden</a> They must instead be replaced with the equivalent character code &. <a href="home.aspx?status=done&itWorked=false">Home & Garden</a> Image alt attributes are mandatory. Because XHTML is designed to be viewed on different types of devices, some of which are not image-capable, alt attributes must be included for all images. <img src="titlebar.gif"> Remember that the img tag must include a closing slash in XHTML! <img src="titlebar.gif" alt="title" /> Scripts and CSS must be escaped. Internal scripts and CSS often include characters like the ampersand and less-than characters. <script language="JavaScript"> document.write('Hello World!'); </script> If you are using internal scripts or CSS, enclose them within the tags <![CDATA[ and ]]>. This will mark them as character data that should not be parsed. If you do not use these tags, characters like & and < will be treated as start-of-character entities (like  ) and tags (like ) respectively. This will cause your page to behave unpredictably, and it may invalidate your code. Additionally, the type attribute is mandatory for scripts. The comment tags that have traditionally been used to hide JavaScript from noncompliant browsers should not be included. The XML standard states that text enclosed in comment tags may be completely excluded from rendered documents, which would lose all script enclosed in the tags. <script type="text/javascript" language="javascript"> /*<![CDATA[*/ document.write('Hello World!'); </script> Also codice_11 is not permitted in XHTML documents. You must used node creation methods such as codice_12 instead. Confusingly, codice_11 will appear to work as expected if the document is incorrectly served with a MIME type of codice_2 (the type for HTML documents), instead of codice_1 (the type for XHTML documents). If the MIME type is codice_2 the document will be parsed as HTML which allows codice_11. Parsing the document as HTML defeats the purpose of writing it in XHTML. Similar changes must be made for internal stylesheets. <style> .SpecialClass { color: #000000; --> </style> The type attribute must be included, and the CDATA tags should be used. <style type="text/css"> /*<![CDATA[*/ .SpecialClass { color: #000000; /*]]>*/ </style> Because scripts and CSS may complicate an XHTML document, it is strongly recommended that they be placed in external .js and .css files, respectively. They can then be linked to from your XHTML document. <script src="myscript.js" type="text/javascript" /> <link href="styles.css" type="text/css" rel="stylesheet" /> Some elements may not be nested. The W3C recommendations state that certain elements may not be contained within others in an XHTML document, even when no XML rules are violated by the inclusion. Elements affected are listed below. When to convert. By now, it probably sounds as though converting an HTML document into XHTML is easy, but tedious. When would you want to convert your existing pages into XHTML? Before deciding to change your entire Web site, consider these questions. If you answered yes to any of the above questions, then you should probably convert your Web site to XHTML. MIME Types. XHTML 1.0 documents should be served with a of codice_1 to Web browsers that can accept this type. XHTML 1.0 may be served with the MIME type codice_2 to clients that cannot accept codice_1 provided that the XHTML complies with the additional constraints in [Appendix C] of the XHTML 1.0 specification. If you cannot configure your Web server to serve documents as different MIME types, you probably should not convert your Web site to XHTML. You should check that your XHTML documents are served correctly to browsers that support codice_1, e.g. Mozilla Firefox. Use 'Page Info' to verify that the type is correct. XHTML 1.1 documents are often not backwards compatible with HTML and should not be served with a MIME type of codice_2. HTML Tidy. When creating HTML, it's very easy to make a mistake by leaving out an end tag or not properly nesting tags. HTML Tidy is a wonderful application that can be used to correct a number of errors with poorly formed HTML documents and convert it into XHTML. Tidy can also format ugly code to be more readable, including code generated by WYSIWYG editors. HTML Tidy can't generate clean code when it encounters problems it isn't sure of how to fix. In these cases, it will generate an error to let you know where the mistake is located in your document. A few examples of problems that HTML Tidy can remedy: HTML Tidy can also be customized at runtime using a wide array of command line arguments. It is capable of indenting code to make it more readable as well as replacing FONT, NOBR, and CENTER tags with style tags and rules using CSS. Tidy can also be taught new tags by declaring them in the configuration file. You can read more about HTML Tidy at the W3C's HTML Tidy site, as well as download the application as a binary or get the source code. There are several sites that offer HTML Tidy as an online service including the W3C and Site Valet. You can also validate your page using the validator available at http://validator.w3.org/. When not to convert. You shouldn't convert your Web pages if they will always be served with a MIME type of codice_2. Make sure you know how to configure your server or server-side script to perform HTTP content negotiation so that XHTML capable browsers receive XHTML marked as codice_1. If you can't set up content negotiation, stick to HTML 4.01. People viewing your Web pages with mainstream browsers will be unable to tell the difference between a valid HTML 4.01 web page and a valid XHTML 1.0 Web page. Make sure the automated tests you run on your site simulate connections from both XHTML-compatible browsers, e.g. Mozilla Firefox, and non–XHTML-compatiable browsers, e.g. Internet Explorer 6.0. This is particularly important if you use Javascript on your Web site. If maintaining two copies of your test suite is too time consuming, don't convert. Bear in mind that valid HTML 4.01 Strict documents generally require less effort to convert to XHTML 1.1 than valid XHTML 1.0 Transitional documents. A valid HTML 4.01 Strict document can only contain elements that are valid in XHTML 1.1, although a few attributes may need changing. XHTML 1.0 Transitional documents on the other hand can contain ten element types and more than a dozen attributes that are not valid in XHTML 1.1. The XHTML 1.0 Transitional codice_25 element alone has six atrributes that are not supported in XHTML 1.1. Don't be pressured into using XHTML by people talking vaguely about "bad practice". Pin them down to what they mean by "bad practice". If they start talking about separation of content and presentation, they have confused the differences between HTML and XHTML with the differences between the Transitional and Strict doctypes. Both XHTML 1.0 Transitional and HTML 4.01 Transitional allow you to mix presentation and content in the same document, i.e. they allow this type of bad practice. Both HTML 4.01 Strict and XHTML 1.0 Strict force you to move the bulk of the presentation (but not all of it) in to CSS or an equivalent language. All four doctypes allow you to use embedded stylesheets, whereas, true separation requires that all CSS and Javascript be moved to external files. XHTML 1.1. XHTML 1.0 is a suitable markup language for most purposes. It provides the option to separate content and presentation, which fits the needs of most Web authors. XHTML 1.1 enforces the separation of content and presentation. All deprecated elements and attributes have been removed. It also removes two attributes that were retained in XHTML 1.0 purely for backwards-compatibility. The codice_6 attribute is replaced by codice_5 and codice_9 is replaced by codice_10. Finally it adds support for ruby text found in East Asian documents. DOCTYPE. The DOCTYPE for XHTML 1.1 is: <!DOCTYPE html PUBLIC "-//W3C//DTD XHTML 1.1//EN" "http://www.w3.org/TR/xhtml11/DTD/xhtml11.dtd"> Modularization. The modularization of XHTML, or XHTML m12n, provides suggestions for customizing XHTML, either by integrating subsets of XHTML into other XML applications or extending the XHTML element set. The framework defines two proceses: The resulting languages, which the W3C calls "XHTML Host Languages", are based on the familiar XHTML structure but specialized for specific purposes. XHTML 1.1 is an example of a host language. It was created by grouping the different elements available to XHTML. XHTML variations, while possible in theory, have not been widely adopted. There is continuing work being done to develop host languages, but their details are beyond the scope of this discussion. Invalid XHTML. XHTML-compliant browsers are allowed to render invalid XHTML documents provided that the documents are well-formed. A simple example is given below: <?xml version="1.0" encoding="UTF-8"?> <!DOCTYPE html PUBLIC "-//W3C//DTD XHTML 1.0 Strict//EN" "http://www.w3.org/TR/xhtml1/DTD/xhtml1-strict.dtd"> <html xmlns="http://www.w3.org/1999/xhtml"> <head> <title>Invalid XHTML</title> </head> <body> <p>This sentence contains a <p>nested paragraph.</p></p> </body> </html> Save the example as codice_30 (the .xhtml extension is important) and open the page with Mozilla Firefox. The page will render even though it is invalid.

Latin/First Declension (L2). =First Declension= The first declension is a group of Latin nouns with the same endings. It is the most regular of the five declensions, in that the twelve endings (6 cases X 2 numbers) are always exactly the same. All you need to do is find the stem. Most first declension nouns are feminine, though a few notable exceptions are masculine. Case Endings. We've already seen the first declension in the charts in the introduction to nouns, but let's take a closer look now. First, the singular: Now, the plural: You will have to memorize these endings; you should do this now. Memorize all twelve endings even though this lesson will only use four (nominative singular, nominative plural, accusative singular, accusative plural). Finding the Stem. How do we determine what the stem (i.e., the part before the ending that remains the same in each case) is and what the ending is? Well, we go back to the dictionary entry that we learn for each word. First, we need to check to see if we have a first declension noun. If the nominative singular ends in "-a" and the genitive singular ends in "-ae", we have a first declension noun. Note how both of those words have the endings that we just said were marks of the first declension nouns. Knowing that "-ae" is the ending for the genitive singular, all we need to do is delete that ending from the word. So, "mensae" loses the "-ae" and becomes "mens-". That is what is known as the stem of the word. You can add the necessary ending for whatever case and number is needed to that stem. So, if you want an ablative plural (case ending "-is"), you would just add that ending to "mens-" to get "mensis". Dictionary Convention. Often, dictionaries will shorten the entry by just printing the genitive singular ending instead of the entire genitive singular. So, "table" and "girl" will appear as: If you see the "-ae" notation in an entry, it is clear that the word is a first declension noun. =Nominative Case= The nominative, or "name" from the Latin "nomen", case is almost always used for the subject of a sentence. For instance, in the sentence "Julia gives" the subject is Julia . She's the one doing the action. So, if the sentence were in Latin, Julia would have to be in the nominative case. Even in a more complex sentence such as "Julia hesitantly decided that it would be a wonderful idea to destroy the table by repeatedly kicking it", we can still find the subject fairly easily. It's still Julia doing the action. Now, we're not even going to try translating that second sentence into Latin quite yet. But, we know most of what we need to know to translate the first one. First, we need to know how to say Julia in Latin. Luckily for us, many women's names were treated as first declension nouns in Latin, so: Julia has a stem of "Juli-" and in the nominative singular (nominative because Julia's the subject, singular because there's only one Julia) the ending is "-a". Therefore, we have "Julia". "Julia gives". We haven't learned any verbs yet and there is a lot to learn about verbs before they can be used well, but for our purpose right now we will just say that dat means "gives". So, if we want to say "Julia gives", the sentence in Latin is: Your first Latin sentence! Pretty simple, huh? Well, let's add on to it. =Accusative Case= The accusative case is often the direct object of a sentence. Take this sentence: "The boy reads a book." We know from our discussion on the nominative case that "the boy" is the subject. "Reads" is our verb, so the boy is reading something. The answer to the question "What is the boy reading?" is the direct object. In this case, the boy is reading a "book". "Book", then, is our direct object and would be in the accusative case in Latin. "Julia gives a table". So, going back to our earlier sentence with Julia, let's pick something for Julia to give. Maybe Julia has a lot of tables and wants to get rid of one. We still don't know who is going to get this table, but we can say "Julia gives a table." Recall that table is: So, knowing that the stem of the word is "mens-" and the accusative single ending is "-am", the table that Julia gives would be mensam. Because the article "a" from "a table" does not translate into Latin, we know all three of the words we need for our short sentence ("dat", "Julia", and "mensam"). But what order do they go in? In all actuality, the words could go in any order and you would be able to translate the sentence, because the case of each of the nouns tells you what its usage is. So, even if we said, "Mensam dat Julia.", you would know that the table cannot be the subject. In actual Latin, the most common sentence structure is nominative-accusative-verb. Compare this to English subject-verb-object. So, if we use the Latin word order, we get: Note that if we decide Julia has to get rid of more than one table, we could say: The Accusative Case is: f.-am -as m.-us -os n.-um -a "The table gives Julia". If we wanted to turn the tables (excuse the pun) and say that the table gives Julia (completely nonsensical, I know, but grammatically correct), we would have to put "Julia" in the accusative ("Juli-am") and "table" in the nominative ("mens-a"). Then, we would get: "Mensa Juliam dat". The table gives Julia. =Vocab= So far we've learned: Let's add a couple more. Make sure you memorize the three above as well as these two following. Wait a minute! Take a look at "poet" again. See something odd? It's not feminine! Poet is one of the few first declension nouns that are masculine. You don't really need to worry about that too much quite yet, because it does not affect how the noun is used. Later, when we get into adjectives and other things, you will need to know what each noun's gender is. For now, use "poeta" like you would any of the others and just keep its masculinity in the back of your mind. =Exercises= With those vocab words and translating "dat" as "gives" and "amat" as "loves", translate the following sentences. Answers =Summary= Congratulations! You now know four common nouns and one name! You know the case endings for every single first declension word, so that even if you come across one that you've never seen before, you will know what case and number it has. You know how to use the nominative as the subject of a sentence, and how to use the accusative as the direct object. In other words, you're ready to start learning some verbs. The next chapter, an introduction to verbs, will teach you the basics of Latin verbs. Lesson 1 - Introduction to Nouns | Latin | Lesson 3 - Introduction to Verbs L2 Word List

Latin/L2/Answers. First Declension (L2) Answers. These are the answers to the exercises in the "First Declension" chapter. Return to First Declension

XML - Managing Data Exchange/XHTML/Answers. XHTML Chapter => XHTML XHTML Exercises => Exercises Answer. 1. Change the HTML document below into an XHTML document conforming to the W3C's "transitional" standard. Validate your page using the validator available at http://validator.w3.org/. <br> Exercise 1. Below is an example of an invalid HTML document that contains a number of deprecated features. It is also badly structured. The document needs converting to XHTML with the content separated from the presentation. Exercise 1.a.. a. Change the HTML document into a valid XHTML 1.0 Transitional document. Validate your page using the validator available at http://validator.w3.org/. The XHTML 1.0 Transitional document: Exercise 1.b.. b. Starting with the model answer for question 1.a. identify all deprecated elements and attributes in the document and replace them with CSS in a linked external stylesheet. See HyperText Markup Language/Tag List for details of deprecated elements. The deprecated features are: A possible stylesheet is: The document with deprecated features removed – note the use of the codice_5 and codice_6 attributes to provide hooks for the CSS selectors. Exercise 1.c.. c. Starting with the model answer for question 1.b. replace all presentational markup with semantic markup and ensure that all inline elements are contained in block-level elements. Change the DOCTYPE to XHTML 1.0 Strict and validate the page with the W3C Markup Validation Service. The presentational elements are: The 'paragraph' beginning with a codice_10 element is a mixture of inline elements and text and so should be contained in a block-level element such as codice_11 (paragraph). The codice_12 (anchor) element is inline and needs enclosing in a block-level element. XHTML Chapter => XHTML XHTML Exercises => Exercises

Calculus/Real numbers. Fields. You are probably already familiar with many different sets of numbers from your past experience. Some of the commonly used sets of numbers are Note that different sets of numbers have different properties. In the set integers for example, any number always has an "additive inverse": for any integer formula_8 , there is another integer formula_9 such that formula_10 This should not be terribly surprising: from basic arithmetic we know that formula_11 . Try to prove to yourself that not all natural numbers have an additive inverse. In mathematics, it is useful to note the important properties of each of these sets of numbers. The rational numbers, which will be of primary concern in constructing the real numbers, have the following properties: As presented above, these may seem quite intimidating. However, these properties are nothing more than basic facts from arithmetic. Any collection of numbers (and operations + and * on those numbers) which satisfies the above properties is called a "field." The properties above are usually called "field axioms." As an exercise, determine if the integers form a field, and if not, which field axiom(s) they violate. Even though the list of field axioms is quite extensive, it does not fully explore the properties of rational numbers. Rational numbers also have an "ordering."' A "total ordering" must satisfy several properties: for any numbers "a," "b," and "c" To familiarize yourself with these properties, try to show that (a) natural numbers, integers and rational numbers are all totally ordered and more generally (b) convince yourself that any collection of rational numbers are totally ordered (note that the integers and natural numbers are both collections of rational numbers). Finally, it is useful to recognize one more characterization of the rational numbers: every rational number has a decimal expansion which is either repeating or terminating. The proof of this fact is omitted, however it follows from the definition of each rational number as a fraction. When performing long division, the remainder at any stage can only take on positive integer values smaller than the denominator, of which there are finitely many. Constructing the Real Numbers. There are two additional tools which are needed for the construction of the real numbers: the upper bound and the least upper bound. Definition A collection of numbers "E" is bounded above if there exists a number "m" such that for all "x" in "E" "x≤m". Any number "m" which satisfies this condition is called an upper bound of the set "E." Definition If a collection of numbers "E" is bounded above with "m" as an upper bound of "E", and all other upper bounds of "E" are bigger than "m", we call "m" the "least upper bound" or "supremum" of "E", denoted by sup "E." Many collections of rational numbers do not have a least upper bound which is also rational, although some do. Suppose the numbers 5 and 10/3 are, together, taken to be "E." The number 5 is not only an upper bound of "E," it is a least upper bound. In general, there are many upper bounds (12, for instance, is an upper bound of the collection above), but there can be at most one least upper bound. Consider the collection of numbers formula_16: You may recognize these decimals as the first few digits of pi. Since each decimal terminates, each number in this collection is a rational number. This collection has infinitely many upper bounds. The number 4, for instance, is an upper bound. There is no least upper bound, at least not in the rational numbers. Try to convince yourself of this fact by attempting to construct such a least upper bound: (a) why does pi not work as a least upper bound (hint: pi does not have a repeating or terminating decimal expansion), (b) what happens if the proposed supremum is equal to pi up to some decimal place, and zeros after (c) if the proposed supremum is bigger than pi, can you find a smaller upper bound which will work? In fact, there are infinitely many collections of rational numbers which do not have a rational least upper bound. We define a real number to be any number that is the least upper bound of some collection of rational numbers. Properties of Real Numbers. The reals are totally ordered. Upper bound axiom The upper bound axiom is necessary for calculus. It is not true for rational numbers. We can also define lower bounds in the same way. Definition A set "E" is bounded below if there exists a real M such that for all "x"∈"E" "x≥M" Any "M" which satisfies this condition is called an lower bound of the set "E" Definition If a set, "E", is bounded below, "M" is an lower bound of "E", and all other lower bounds of "E" are less than "M", we call "M" the "greatest lower bound" or "inifimum" of "E", denoted by inf "E" The supremum and infimum of finite sets are the same as their maximum and minimum. Theorem "Proof:" formula_17 formula_18 formula_19 Theorem: "(The triangle inequality)" "Proof" by considering cases If a≤b≤c then "|a-c|+|c-b|" = "(c-a)+(c-b)" = "2(c-b)+(b-a)≥b-a" = "|b-a|" "Exercise:" Prove the other five cases. This theorem is a special case of the triangle inequality theorem from geometry: The sum of two sides of a triangle is greater than or equal to the third side. It is useful whenever we need to manipulate inequalities and absolute values.

Electronics/Preface. Importance of Electronics. Electronics is the study and use of devices that control the flow of electrons (or other charged particles). These devices can be used to process information or perform tasks using electromagnetic power. Electronic circuits can be found in numerous household products, including such items as telephones, computers, and CD players. Electronic devices have also allowed greatly increased precision in scientific measurements. Interest in the field of electronics increased around 1900 and the advent of radio. Interest reached an all-time high in the 1940s, 50s, 60s, with the invention of transistor radios, the launch of Sputnik, and the science and math educational push to win the space race. Interest in electronics as a hobby in the 1970s led to the advent of the personal computer (PC). Electronics have since seen a decline in hobbyist interest. Electronics is now generally studied as part of a college-level program in electrical engineering. This book is an attempt at reviving the hobbyist mentality that made electronics so big in the first place, by making electronics concepts more accessible and giving practical knowledge, as well as providing technical information for the student.

Latin/Introduction to Verbs (L3). =Introduction to Verbs= Verbs in Latin, like verbs in English, are the "action words". "Go", "throw", "kick", "fly", "think", and "be" are all examples of verbs. Verbs in Latin. Just like nouns, we can tell a lot about a Latin verb by looking at its ending. This is because verbs conjugate (remember, nouns decline). There are four groupings of verbs, called "conjugations". Like the declensions for nouns, each verb conjugation contains a group of verbs with similar endings. We saw that nouns have three incidents. Verbs, however, have five. As a result, there is a vast number of forms that each verb can take. We'll start slow, of course. First, before we get into talking about verbs themselves, we should talk about what those five incidents are. Mood. Verbs have three different moods - the indicative, the imperative, and the subjunctive. The indicative, from the word "indicate", is the mood that shows action happening. When Bob kicks a ball, "kicks" is a verb in the indicative because it tells us that something is actually happening. The imperative is a command. If I say to Bob, "Kick the ball", "kick" is in the imperative - I'm telling Bob to kick the ball, the ball isn't actually being kicked yet. Finally, Latin verbs have a subjunctive mood. This usually represents hypothetical action. "Bob might kick the ball" uses "kick" in the subjunctive mood. We'll learn more about the indicative mood later in this lesson. Voice. Just as in English, there are two voices that verbs can take. The active voice is used when the subject of a sentence (the nominative case noun) is doing the action. Such as "Bob kicked the ball" - Bob is the subject and he is doing the kicking. But when we say "The ball was kicked by Bob", the ball is the subject of the sentence. The ball isn't doing anything though! In this case, the verb "was kicked" is passive - the subject has something done to it. Note that in English we need to use the auxiliary (helping) verb "was" along with our main verb "kick" to show the passive voice. In Latin, a form of "kick" is used that shows passive voice without needing an auxiliary. We'll start off by focusing on the active voice. Tense. We're familiar with tense in English - we can tell the difference between "Bob kicks the ball", "Bob kicked the ball", and "Bob will kick the ball". Those are just a few of the tenses that we use in English - the present, the past, and the future. Latin has a total of six tenses. We'll go into each of the tenses later (starting with the present in this lesson), but for now just know that there are six. Person. "I kick the ball". "You kick the ball". "Bob kicks the ball". What's the difference in these three sentences? The person of the verb. We all remember from our English classes the different points-of-view that can be used in writing. First person is the use of words that refer to the person speaking (I, me, we, us). Second person refers to the person being spoken to (you). Third person talks about someone else entirely (Julia, Bob, the girl, the poet). Note in the three sentences above the first and second person in English use "kick" while third person uses "kicks" - this is a vestige of the Latin conjugation system. In Latin, each of those three sentences would use a different form of "kick". The person of a verb is determined by the subject of a sentence. Number. Just like nouns, verbs have number. They can be singular or plural. The number of the verb is determined by the number of the subject of the sentence - if you have a singular subject, you need to use a singular verb. Multiple subjects require a plural verb. In English, we really only see this in the third person ("Bob kicks the ball" vs. "Bob and Julia kick the ball"). However, Latin has two different forms for each person. =First Conjugation= The first conjugation is the most regular of the four groupings of verbs. Both of the verbs we used in our last lesson belong to the first conjugation. Can you tell from the sentences what person and number those verbs were? How about tense, voice, and mood? If you decided that they were third-person singular, you're right. Bonus points for present active indicative. Let's stay with the present active indicative for right now, and take a look at the different endings for person and number. Present Active Indicative. Present active indicatives are forms of verbs that show a subject doing an action now. They can be translated into English either as simple presents (something happens) or as progressive presents (something is happening). Note that "you" in English can be singular or plural, so we will distinguish it here as follows: "you" - singular, "y'all" - plural. The "dat" should look familiar because we used it in the last lesson. We also used "amat". From looking at the chart above, can you guess what the other forms of "amat" are? If that's what you came up with, you did well. Dictionary Entries. The dictionary entry for a Latin verb has four "principle parts". Once we get into the later lessons, you will need all four parts to be able to fully conjugate (i.e., list all the forms of) a verb. For now, only the first two are important. They are the first-person singular present active indicative and the present active infinitive. We'll talk more about infinitives, and why they don't have person or number, later. The dictionary entries for the two verbs we've used so far are: The "--"s will be filled in later when we learn the rest of the principle parts. Stem. Notice anything similar about the two entries? The second part (the infinitive) of each ends in "-are". Any time that you see that ending for the infinitive, you will know that the verb is first conjugation. To find the stem of the verb, you simply slash off the "-re" from the infinitive. This holds true for all four conjugations. What are the stems of our two verbs? "dare" --> "da-"<br> "amare" --> "ama-" Personal Endings. What can we do with this stems? Well, to get each of the person/number combinations we simply add a "personal ending" to the verb's stem. The endings are: One issue you might notice is in the first-person singular. Take a look at the stem of "amare". It is "ama-". If we were to add the first-person singular ending to that, we would get "ama-" + "-o" = "amao". But we know from our dictionary entry that the first-person singular is "amo". What's up? In Latin, sometimes certain vowels get "consumed" by other vowels that are stronger. This is what has happened in this case: the "-o" is a strong enough sound that it simply obliterates the "a" in the stem. Because all first conjugation verbs have stems that end in "a" (remember that for a verb to be first conjugation, the infinitive must end in "-are" and so cutting the "-re" will always leave "-a-"), every first conjugation verb will exhibit this change. Missing Subjects. Often in Latin, you will see sentences that have no noun in the nominative case. Don't panic; there is still a subject for that sentence. Let's take a look at an example. The sentence "Sapientiam amo" has an accusative ("sapientiam") and a verb ("amo"), but no nominative. You can find the subject by looking at the verb - "amo" is the first-person singular. There is really only one option for who subject is. It has to be "I" because "I" is the only subject that will use a first-person singular verb ("we" uses a first-person plural"). Therefore, the sentence translates to "I love wisdom". You can use the same logic easily with second-person verbs as well. Third-person is a bit trickier because it has more options, but the singular is "he", "she", or "it" and the plural is "they". Practice. Let's learn a new verb now. How about this one: What is the stem? What is the third-person singular? The second-person plural? =Vocab= In addition to the three verbs you learned earlier in this lesson, add this one to your memory: Let's add another couple of nouns as well. How about the proper nouns "Rome" and "Italy"? =Exercises= Answers =Summary= You now know what a verb is, what conjugation means, and how to conjugate the present active indicative forms of first conjugation verbs. Your noun vocab is expanding, and you've learned your first four verbs. You can translate sentences with subjects, direct objects, and present tense verbs. You can even fill in gaps when the subject of a sentence is missing by looking at the verb. In the next lesson, we'll learn about the dative case and indirect objects. We'll finally discover to whom Julia is giving those tables! Lesson 2 - First Declension | Latin | Lesson 4 - Dative Case and Indirect Object L3 Word List

Latin/L3/Answers. Introduction to Verbs (L3) Answers. These are the answers to the exercises in the "Introduction to Verbs" chapter. Return to Introduction to Verbs

Scouting/BSA/Orienteering Merit Badge. Using a Shadow to Find Directions. 1. Find a straight stick, preferably one about a meter long. 2. Find an area to place your stick. This area ahould be as open as possible, free of brush or trees. It should also be as level as possible. 3. Place the stick in the ground sticking up. 4. Mark the tip of the shadow cast by the stick. 5. Wait for fifteen or twenty minutes until the shadow has moved noticeably. 6. Mark the tip of the shadow again. 7. Draw a line between these two marks. The first mark is west and the second is east. 8. If you stand with the first mark directly to the left of the second mark you will be facing north.

Outdoor Survival. Related Wikibooks. __NOEDITSECTION__

Electronics/Power Conversion. AAC. Analog-to-Analog Conversion is the conversion of a continuous electrical signal from one form to another. An example would be converting between the 60 Hz current of the US into the 50 Hz current of Europe. ADC. Analog to digital Used in power supplies involves a transformer to change the voltage and diodes to rectify the voltage. Analog to Digital Conversion is the process of making a computerized representation of some physical attribute in the real world. Turning a voltage into binary numbers is the handiest way to do this, the voltage could be manipulated by a microphone transducer, a potentiometer wiper, or the conductivity of prevarication. DAC. Digital to analog DDC. digital to digital

Electronics/Timeline. < Electronics/Expanded Edition Different ways of looking at electronics.

Electronics/Electromagnetic Waves. < Electronics/Timeline Originally people thought that radio waves propagate through the ether. The ether was disproved by showing that it did not flow. EM waves are photons that travel through free space.

XML - Managing Data Exchange/Recursive relationships/Answers. Exercises 2a. Answer: File relayXML.xml File relayXSD.xsd File relayXSL.xsl Answer: File insuranceXML.xml File insuranceXSD.xsd File insuranceXSL.xsl Answer: File appetizerXML.xml File appetizerXSD.xsd

HyperText Markup Language/Images. Let us start with a quick minimum example: <img src="turtleneck-sweater.jpeg" /> And let us also have a look at more extended example: <img src="turtleneck-sweater.jpeg" alt="Photograph of a Turtleneck Sweater" title="Photograph of a Turtleneck Sweater" /> Images are normally stored in external files. To place an image into an HTML source, use the codice_1 tag with the codice_2 attribute containing the URL of the image file. To support browsers that cannot render images, you can provide the codice_3 attribute with a textual description of the image. To provide a tooltip to the image, use the codice_4 attribute. The space before the codice_5 in the examples is there on purpose. Some older browsers behave strangely if the space is omitted. Placement. Per default, images are placed "inline" to their surroundings. To place the image as a "block" or "float" instead, you can use Cascading Style Sheets. Alternative text and tooltip. The HTML specification requires that all images have an codice_3 attribute. This is commonly known as "alt text". Images can be split in to two categories: Decorative images should have empty alt text, i.e. codice_7. Images used as bullets may use an asterisk, codice_8. All other images should have meaningful alt text. Alt text is not a description of the image, use the codice_4 attribute for short descriptions. Alt text is something that will be read instead of the image. For example, <img src="company_logo.png" alt="???"> makes the best widgets in the world. The alt text should be the company's name not the ever popular 'Our logo', which would give the sentence 'Our logo makes the best widgets in the world.' when read in a text only browser. The codice_3 attribute stands for "alternate" which screenreaders for the blind, as well as non-graphic-capable browsers (such as Lynx) may use to better enable its user to understand the purpose of the image. Fully and accurately using alt text is essential for the following reasons: Tooltips. Some older browsers incorrectly use the codice_3 attribute' tag to produce image tooltips. However, the codice_4 attribute should actually be used for this. Width and height. To have an image appear smaller than it is in the external file, use the attributes codice_13 and codice_14. When only one or the other is specified, images will scale to keep the same overall ratio. Format. Most images for the web will be made using a limited number of formats. Below is a simple list of commonly supported formats.

Electronics/Expanded Edition. Chapter 13: Semiconductors. Semiconductors are often made of silicon. Chapter 15: Robotics. An excellent source of info is the various Robotics Societies, including Seattle Robotics Society and Dallas Personal Robotics Group. There is a lot of info on these groups websites, especially in the archived pages of The Encoder, published by SRS. To begin, an excellent source would be the Wikipedia pages on Robotics (w:Robotics) or the Robotics Wikibook. Chapter 16: Spacecraft. Telemetry - Measurement at a distance. Critical in the success of early spacecraft. This technology allowed an incredible pyramid of resources to be applied to monitoring and safely operating the machine and human occupied spacecraft of the pre-space era of human development. Modern development efforts are focusing on a new generation of space based grid technologies that are hopefully sustainable independently from the mother planet with only occasional unnoticed intervention from the Lord Almighty but which when prematurely, suddenly or unexpectedly exposed to an adverse probability surge from outside the locally anticipated systems or universes can survive long enough to request and receive assistance. Telecommunications []

Electronics/Electrolysis. Electrolysis involves an electric field and the motion of electrons through a liquid. GCSE Science/Electricity Electrolysis is the decomposition of certain types of substance using electricity.The types of substance that can be split are ionic substances.This just means that they are made of charged ions rather than neutral atoms.{Remember that an ion is just an atom that has either a positive or negative charge}. An example of an ionic substance is common table salt sodium chloride. The sodium atom has a positive charge, the chlorine atom has a negtive charge. It is usually written as Na+Cl-. Q1) Check in a periodic table, what is the symbol for sodium Na or Cl? As you may already know if you've studied the Metals module, a salt is any substance made by combining an acid with an alkali. Acids, alkalis, and therefore all salts are ionic. Q2) Which of the following substances can be broken up by electrcity Sodium chloride, iron sulphate, copper nitrate? Basic Experimental Setup. Most ionic compounds are not liquid at room temperature. This is a problem because the ions need to be able to move for the electric current to be able to flow. This can be achieved by melting. Look at the electrical setup shown on the right. The electrodes are just tqwo carbon rods connected to a battery. The one connected to the positive electrode is called the anode. The one connected to the negative electrode is called the cathode. Consider for example the compound lead bromide. This compound is a solid at room temperature but can be molten over a Bunsen flame. So what you would do is put some lead bromide into a beaker. Put the beaker on a tripod over a bunsen flame. Melt the lead bromide, then put in the electrodes and turn the power supply on at a setting of say 2V. What you would see happening is the anode, is a silvery coating of pure lead forming, and bromine forming at the cathode. The current would continue to flow until all the lead bromide was turned into lead and bromine. Q3) It takes energy to split uo a compound like lead bromide. Where does this energy come from ? Q4) Predict what products you would get at the anode and cathode if copper chloride was the electrolyte. What happens at the anode. The anode is the positive electrode. It attracts negatively charged ions, because unlike charges attract. The bromine ions move through the melt until they reach the anode. Once they get there, they give up their two extra electrons to become bromine atoms. Br2- --> Br + 2e- The electrons flow up the anode to the positive terminal of the battery. What happens at the cathode. The cathode is the negative electrode. It attracts the positively charged ions. Metal ions are always positive and so the lead ions flow through the melt until they get to the cathode. Once they get there they take two electrons each from the cathode. The electrons flow down the cathode from the negatively terminal of the battery and onto the lead ions. Pb2+ + 2e- --> Pb Q5) Solid ionic substances do not conduct electricity and are not split up by it. Why do you think that is? Quantity calculations (higher tier only). In the experiment with lead bromide, you saw that lead was deposited at the cathode. If you actually do the experiment you will see that the lead coats the cathode. In this section we will look at how much metal will coat a cathode in a given time. A scientist performed the following experiment. His results were: You can see from the results that the total amount of copper deposited depends on both the current and the time it flows. This is because the number of copper atoms that can be made from ions depends on the total amount of charge that flows. The unit of charge is the coulomb. One coulomb is the amount of charge when one Amp flows for one second (see GCSE Science/Measuring electricity Q6)Look at the results table above. How much copper is deposited when 1A flows for 3000 seconds? Q7)How much copper do you predict would be doposited if 1A were to flow for 6000 seconds. Q8)What about if 2A were to flow for 12000 seconds ? Electrolysis of Aqueos solutions (Advanced). "Before studying this section check with your teacher to see if you need to". Earlier on in this module you've learned that ions must be able to move in order for electrolysis to work. If the ions are held rigid {such as in a solid}, they can't move and no electricity will flow. We've looked at how the freeing up of ions can occur by melting the electrolyte. Another way to achieve this is by dissolving th electrolyte in water. The trouble with this method is, there will be more than one type of ion present. water partially splits up into ions {this is why it's such a good solvent for ionic compounds}. It splits into hydrogen ions and hydroxide ions. H2O --> H+ +OH- So at the cathode there will be two ions present. The Metal ion and the hydrogen ion from the water. Which element is actually produced at the cathode depends on how reactive the metal is. If the metal is very reactive, such as potassium, or sodium, then hydrogen will be produced. If the metal is unreactive such as silver, the metal will be produced.{I hope you learned your displacement reactions in lower school}. To work out which ion wins, the metal or the hydrogen, compare their reactivities in a the reactivity series. The one that is most reactive, will not be produced at the cathode. A similar situation occurs at the anode. Hydroxide ions {from the water} are usually discharged at the anode ultimately producing oxygen. OH- --> OH + e- 4OH --> 2H2O + O2. Q9) Sodium chloride is dissolved in water and subjected to electrolysis. Explain what you see at each of the electrodes. Answers | «Advanced static electricity | Circuits»

Statistics/Displaying Data.      A single statistic tells only part of a dataset’s story. The mean is one perspective; the median yet another. And when we explore relationships between multiple variables, even more statistics arise. The coefficient estimates in a regression model, the Cochran-Maentel-Haenszel test statistic in partial contingency tables; a multitude of statistics are available to summarize and test data. <br><br>      But our ultimate goal in statistics is not to summarize the data, it is to fully understand their complex relationships. A well designed statistical graphic helps us explore, and perhaps understand, these relationships. <br><br>      This section will help you let the data speak, so that the world may know its story. <br> Types of Data Display External Links Statistics | » Bar Charts

Lucid Dreaming/Reality Checks/Light switches. Presentation. With the "light switches" reality check, you check if light switches work normally. This is one of the less reliable reality checks, but for some people it works perfectly.

Electronics/Cables. Everything past this point is notes. Gage of wire. The larger the gage the smaller the wire if you are using AWG. Metric sizes go the other way. Cables. Component video and audio. Coaxial. Coaxial cable is an electrical cable consisting of a round, insulated conducting wire surrounded by a round, conducting sheath, usually surrounded by a final insulating layer. The cable is designed to carry a high-frequency or broadband signal, usually at radio frequencies. Sometimes DC power (called bias) is added to the signal to supply the equipment at the other end, such in direct broadcast satellite receivers. Because the electromagnetic field carrying the signal exists (ideally) only in the space between the inner and outer conductors, it cannot interfere with or suffer interference from external electromagnetic fields. Coaxial cables may be rigid or flexible. Rigid types have a solid sheath, while flexible types have a braided sheath, both usually of copper. The inner insulator, also called the dielectric, has a significant effect on the cable's properties, such as its characteristic impedance and its attenuation. The dielectric may be solid or perforated with air spaces. Coaxial cables are usually terminated with RF connectors. Electromagnetics. Open wire transmission lines have the property that the electromagnetic wave propagating down the line extends into the space surrounding the parallel wires. These lines are low loss, however they have undesirable characteristics. They cannot be bent, twisted or otherwised shaped without changing their characteristing impedance. They also cannot be run along or attached to anything conductive, as the extended fields will induce currents in the nearby conductors causing unwanted radiation and detuning of the line. Coaxial lines solve this problem by confining the electromagnetic wave to the area inside the cable, between the center conductor and the shield. The line itself forms a coaxial waveguide, and the transmission of energy in the line occurs totally through the wave that propagates inside the cable between the conductors. Coaxial lines can therefore be bent and twisted without negative effects, and they can be strapped to conductive supports without inducing unwanted currents on them. Coaxial lines are filled with a dielectric material that maintains the spacing between the center conductor and shield. Unfortunately, all dielectrics have loss associated with them, which causes most coaxial lines to be lossier than open wire lines. Standard cable types. Most coaxial cables have a characteristic impedance of either 50 or 75 ohms. The RF industry uses standard type-names for coaxial cables. The U.S. military uses the RG-# or RG-#/U format (probably for "radio grade, universal", but other interpretations exist). For example: ("φ" = diameter, ΩΩ = ohms) Uses of coaxial cable. Short coaxial cables are commonly used to connect home video equipment, or in ham radio setups. Long distance coaxial cable is used to connect radio networks and television networks, though this has largely been superseded by other more high-tech methods (fibre optics, T1/E1, satellite). In broadcasting and other forms of radio communication, hard line is a very heavy-duty coaxial cable, where the outside shielding is a rigid or semi-rigid pipe, rather than flexible and braided wire. Hard line is very thick, typically at least a half inch or 13mm and up to several times that, and has low loss even at high power. It is almost always used in the connection between a transmitter on the ground and the antenna or aerial on the tower. Hard lines are often made to be pressurised with nitrogen or desiccated air, which provide an excellent dielectric even at the high temperatures generated by thousands of watts of RF energy, especially during intense summer heat and sunshine. Physical separation between the inner conductor and outer shielding is maintained by spacers, usually made out of tough solid plastics like nylon. Triaxial cable also exists, in which a third layer of insulation and sheathing is included. This allows a nearly perfect signal which is both shielded and balanced/differential to pass through. Multi-conductor coaxial cable is also used sometimes. Biaxial cable or biax is a figure-8 configuration of two 50 ohm coaxial cables, used in some proprietry computer networks. Timeline. Copper Fiberoptics Cat 5 Cable Configurations. Twisted Pair. Twisted pair cabling is a common form of wiring where two conductors are wound around each other for the purposes of cancelling out electromagnetic interference known as crosstalk. The number of twists per meter make up part of the specification for a given type of cable. The greater the number of twists, the more crosstalk is reduced. Shielded twisted pair (STP) has an outer conductive casing similar to coaxial cable and theoretically offers the best protection from interference. It was commonly used for token ring networks. UTP or unshielded twisted pair is not surrounded by shielding. It is the primary wire type for telephone usage and is very common for computer networking. UTP is standardized into various categories by number, which indicate signal integrity attributes. Category 5 cable is commonly used for Ethernet with 10BASE-T or 100BASE-TX. In telephone applications, UTP is often grouped into sets of 25 pairs according to a standard 25 pair color code originally developed by AT&T. A typical subset of these colors (white/blue, blue/white, white/orange, orange/white) shows up in most UTP cables. Audio cables that are twisted are often referred to as balanced. Connectors. Registered Jack (RJ) - Untwisted and Twisted pair GCSE Science/Electricity Everyone knows mains is potentially dangerous. There are however safety features included in plugs. This module looks at how to correctly wire a plug. The three pin plug. Nowadays most appliances are sold with moulded plugs already fitted. Nevertheless it is still important to understand the correct wiring of a plug because enough of the old plugs still exist. You are very likely to need to change a plug at some time in your life. In the UK mains electricity is 230V. if you were to touch the live wire a current would flow through your body to the ground. This current may be enough to kill you. The cable from the appliance usually consist of three wires. The wires are made of copper surrounded by a plastic sheath. The sheath is made of plastic and is coloured: The three wires are covered by an outer sheath made of plastic. Q1)Use you knowledge of insulators and conductors to explain The plug has the following features: Q2) Why are the pins made of brass and why is the case plastic? The live and neutral wires. They carry the current around the circuit. Mains current is AC this means that it is going backward and forwards in cycles. (clockwise and anticlockwise around the circuit). The frequency of the cycle is 50 times a second (50 Hertz) This cycling of current is achieved by varying the voltage on the live wire from +230V to -230V and back again. The voltage on the neutral wire does not vary. It stays close to zero ( hence the term neutral). In contrast the voltage of a battery does not cycle. It says constant. This is known as D.C. (direct current) The live and neutral pins of a plug often have the half nearest the case surrounded by a plastic shroud. This is to prevent your fingers from touching them if you curl them around the plug as you push the plug into the socket. Q3)Why is only the part nearest the cases sheathed. Why not the whole pin? The earth wire. This wire is there to protect you . Many appliances have metal cases. For example kettles, toasters, dishwashers and washing machines. If the live wire were to become loose inside the appliance and touch the case the whole case would become live. If you were then to touch it a current would flow through you to the earth. The earth wire is just a wire connected to the case of the appliance. It goes down the flex into the socket. In side the wiring of your house it travels down to the earth. It is often connected to a metal plumbing pipe. If the live wire were to touch the case a huge current would flow through the earth wire. This would probably blow the fuse and break the circuit {see next section} but even if the fuse doesn't blow the current would still prefer to flow through a nice low resistance wire than a high resistance human body. The earth pin on a plug is longer than the live and neutral pins. This ensures that the earth pin always connects with the socket first. The fuse. A fuse is simply a very thin wire. The wire has a low melting point. As current flows through the wire it heats up. If too large a current flows it melts. This breaks the circuit. Fuses are used to protect appliance. If too large a current flows through an appliance it may damage it. It may even cause a fire! Fuses are rated according to how much current they can carry before melting. In plugs fuses are usually 3A , 5A, or 13A. The correct fuse for an appliance is one that is just above the normal working current for that appliance. Q4) A table lamp usually carries a current of 0.5 A what fuse should be put in the plug 3 A, 5A, or 13A? Q5) An iron usually carries a current of 5.2 A what fuse should be put in the plug 3 A, 5 A, or 13 A? Q6) A kettle is protected by an earth wire and a 13 A fuse. The live wire comes loose and touches the side of the kettle. The fuse blows. Explain why. The unit is designed to break the circuit if the current spikes above 13A.

Statistics/Displaying Data/Bar Charts. Bar Charts. The Bar Chart (or Bar Graph) is one of the most common ways of displaying catagorical/qualitative data. Bar Graphs consist of 2 variables, one response (sometimes called "dependent") and one predictor (sometimes called "independent"), arranged on the horizontal and vertical axis of a graph. The relationship of the predictor and response variables is shown by a mark of some sort (usually a rectangular box) from one variable's value to the other's. To demonstrate we will use the following data(tbl. 3.1.1) representing a hypothetical relationship between a qualitative predictor variable, "Graph Type", and a quantitative response variable, "Votes".<br> tbl. 3.1.1 - Favourite Graphs From this data we can now construct an appropriate graphical representation which, in this case will be a Bar Chart. The graph may be orientated in several ways, of which the vertical chart (fig. 3.1.1) is most common, with the horizontal chart(fig. 3.1.2) also being used often<br><br> fig. 3.1.1 - vertical chart<br> <br> fig. 3.1.2 - horizontal chart<br> <br> Take note that the height and width of the bars, in the vertical and horizontal Charts, respectfully, are equal to the response variable's corresponding value - "Bar Chart" bar equals the number of votes that the Bar Chart type received in tbl. 3.1.1 Also take note that there is a pronounced amount of space between the individual bars in each of the graphs, this is important in that it help differentiate the Bar Chart graph type from the Histogram graph type discussed in a later section.

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Electronics/Op-Amps. Op-Amp (operational amplifier). Op-amp stands for operational amplifier. It is available in IC (Integrated Circuit) chip. It is an electronic circuit of many electronic components already connected and packaged inside a chip of many pins for external connection. Originally, op-amps were so named because they were used to model the basic mathematical operations of addition, subtraction, integration, differentiation, etc. in electronic analog computers. In this sense a true operational amplifier is an ideal circuit element. 741 Op-amp. Symbol. The 741 op-amp has a symbol as shown. Its terminals are: Operation. The 741 op-amp can be thought as Universal Voltage Difference Amplifier. The main function of the 741 op-amp is to amplify the difference of two input voltages which can be expressed mathematically like below From It can be shown that an op-amp can function as a noninverting voltage amplifier or inverting voltage amplifier. Also, it can be shown that an op-amp can function as a voltage comparator Ideal Op-amps. The output voltage is the difference between the + and - inputs multiplied by the open-loop gain The ideal op-amp has The model for the op-amp is shown in Figure 1. Where V+ − V− is equal to vd; Rin is the input impedance; Rout is the output impedance; Avo is the open loop gain; and Rs is the source impedance. Using voltage divider rule, we can determine the voltage vd. Since the dependent voltage source amplifies the voltage vd by Avo. We can work out the output voltage across RL once again using voltage divider rule. Substituting equation 1 into 2. Now if the properties of an Ideal op-amp are applied. The ideal properties of an op-amp are infinite input impedance and zero output impedance. Since Rin », much greater than, Rs, Rin/(Rin+Rs) formula_14 1. Which is basically the definition of an op-amp. But if the input impedance was not infinite and the output impedance nonzero then this would not be so. Real Op-amps. Real op-amps are (normally) built as an , but occasionally with discrete transistors or vacuum tubes. A Real op-amp is an approximation of an Ideal op-amp. A real op-amp does not have infinite open loop gain, infinite input impedance nor zero output impedance. Real op-amps also create noise in the circuit, have an offset voltage, thermal drift and finite bandwidth. An offset voltage Thermal drift Finite bandwidth Modern integrated circuit MOSFET op-amps approximate closer and closer to these ideals in limited-bandwidth, large-signal applications at room temperature. When the approximation is reasonably close, we go ahead and call the practical device an 'op-amp', forget its limitations and use the thinking and formulae given in this article. DC Behaviour. Open-loop gain is defined as the amplification from input to output without any feedback applied. For most practical calculations, the open-loop gain is assumed to be infinite; in reality, however, it is limited by the amount of voltage applied to power the operational amplifier, i.e. Vs+ and Vs- in the above diagram. Typical devices exhibit open loop DC gain ranging from 100,000 to over 1 million. This allows the gain in the application to be set simply and exactly by using negative feedback. Of course theory and practice differ, since op-amps have limits that the designer must keep in mind and sometimes work around. AC Behaviour. The op-amp gain calculated at DC does not apply at higher frequencies. This effect is due to limitations within the op-amp itself, such as its finite bandwidth, and to the AC characteristics of the circuit in which it is placed. The best known stumbling-block in designing with op-amps is the tendency for the device to resonate at high frequencies, where negative feedback changes to positive feedback due to parasitic lowpasses. Typical low cost, general purpose op-amps exhibit a gain bandwidth product of a few MHz. Specialty and high speed op-amps can achieve gain bandwidth products of 100s of MHz. Quick Design Process. Imagine we have a (weak) input signal, and we want to amplify it to generate a strong output signal; or we have several different voltages, and we want to add them together. For these purposes, we need op-amp amplifying and summing circuits. It seems that there is nothing new to say about these bare (usually single op-amp) circuits. Only, the innovator Dieter Knollman has suggested a new simpler design procedure in an EDN's article . He has also placed his work on the web (see also a circuit story written after Dieter's work). We assume: Design procedure: Therefore, choose Some op-amp circuits need a resistor to ground from the op-amp's inverting input. Others need a resistor to ground on the noninverting input. The sign of the ground gain determines where to place the ground resistor. If the desired gains add up to one, a resistor to ground is unnecessary. You will see below some examples that illustrate using the quick design procedure. First of all, answer the question, "What's the highest and lowest output voltage you want?" Make sure the power terminals of the op-amp are at least that high and low. (If you want it to swing from +5 V to -5 V, then an op-amp connected to 0 V and +10 V won't work). The generic op-amp has two inputs and one output. (Some are made with floating, differential outputs.) The output voltage is a multiple of the difference between the two inputs: "floating outputs? Are you perhaps misunderstanding "fully differential amplifier" ?" G is the open-loop gain of the op-amp. The inputs are assumed to have very high impedance; negligible current will flow into or out of the inputs. Op-amp outputs have very low source impedance. If an Operational Amplifier is used with positive feedback it can act as an oscillator. Other Notation. The power supply pins (VS+ and VS−) can be labeled many different ways. For FET based op-amps, the positive, common drain supply is labeled VDD and the negative, common source supply is labeled VSS. For BJT based op-amps, the VS+ pin becomes VCC and VS− becomes VEE. They are also sometimes labeled VCC+ and VCC−, or even V+ and V−, in which case the inputs would be labeled differently. The function remains the same. Fully Differential Amplifier. Basic Applications. "... insert drawings of basic FDA amplifying circuits here ..." Fully Differential Amplifier versus Instrumentation Amplifier. An instrumentation amplifier has extremely high input impedance-- much higher than the input impedance of a fully differential amplifier (once all the feedback resistors are in place). So an instrumentation amplifier is better for measuring voltage inputs with unknown (and possibly time-varying) output resistance. A fully differential amplifier is better than an instrumentation amplifier for precisely generating differential output voltages, with good rejection of differential noise on the input, output, and power lines. Applying a Quick Design Procedure. Can we apply Dieter's procedure to a fully differential amplifier? We can, if we revise it a little: Here is the FDA quick design procedure: Example 7 (a classic design procedure): Example 8 (a quick design procedure):

UK Constitution and Government/Sovereign. The Sovereign. The role of head of state in the United Kingdom is held by the Sovereign; the present Sovereign is King Charles III. Succession to the throne. As a hereditary monarchy, the rules for succession to the throne are established by common law, as modified by statute. In accordance with the "Act of Settlement" (1701), on the death of the Sovereign he or she used to be succeeded by his or her "heir of the body"; this operated in accordance with the principle of "male-preference primogeniture". If the Sovereign had only one child, that child succeeds. If there are more than one children, then the order of succession was determined first by sex, and then by age. This was repealed by the Succession to the Crown Act 2013. If the Sovereign dies childless ("without issue") then the order of succession is applied to their siblings: the oldest surviving sibling then succeeds. If the Sovereign's siblings have died before he or she died, then the order of succession works through the children of the next oldest deceased brother, and so on. Only legitimate children are able to succeed. The "Royal Marriages Act 1772" operates to restrict the capacity for a potential heir to marry without the Sovereign's approval: all descendants of King George II, other than women who have married into foreign families, are required to obtain the Sovereign's consent before marrying, unless they can otherwise obtain approval from both Houses of Parliament. The "Bill of Rights" (1689) and "Act of Settlement" require all heirs to be descendants of Sophia, Electress of Hanover (d. 1714), and impose further requirements that an heir be a Protestant, that they may never have married a Roman Catholic, and that they be in full communion with the Church of England. Heirs not meeting these conditions are skipped over as if "naturally dead". The role of the Sovereign beyond the United Kingdom. As Sovereign in right of the United Kingdom, the Sovereign is also head of state in the "Crown dependencies" of Jersey, Guernsey (and its dependencies), and the Isle of Man. While the external relations of these islands is dealt with by the United Kingdom, however, they do not form part of the United Kingdom itself, and have their own constitutional arrangements. Similarly, the United Kingdom has sovereignty over various territories around the world, known as the "British overseas territories". As such, the Sovereign is also head of state in these territories, although again these do not form part of the United Kingdom itself, and have their own constitutional arrangements. The British Sovereign is also the Sovereign of certain other "Commonwealth Realms": Antigua and Barbuda, Australia, the Bahamas, Barbados, Belize, Canada, Grenada, Jamaica, New Zealand, Papua New Guinea, Saint Kitts and Nevis, Saint Lucia, Saint Vincent and the Grenadines, the Solomon Islands and Tuvalu. Each of these nations is a separate monarchy; the Sovereign therefore holds sixteen different crowns. In each nation, the Sovereign is represented by a Governor-General, who generally stands in relation to the local government in the same relation as the Sovereign does to the British government. Finally, the Sovereign has the title "Head of the Commonwealth". The Commonwealth is a body of nations mostly made up of former colonial dependencies of the United Kingdom. The role of Head of the Commonwealth is a personal role of the present King, Charles III, and is not formally attached to the monarchy itself (although the present King's mother, Queen Elizabeth II, also held the title). The role is purely a ceremonial one. Royal Family. While the members of the Sovereign's family do not have any role in government, they do exercise ceremonial functions on his or her behalf. A male Sovereign has the title "King", while a female Sovereign is the "Queen". The wife of a King is also known as a Queen; however, the husband of a female Sovereign has no specific title. By convention, the Sovereign's eldest son is created "Prince of Wales" and "Earl of Chester" while still a boy; he also automatically gains the title of "Duke of Cornwall". Also by convention, the Sovereign's sons receive a peerage either upon reaching the age of twenty-one, or upon marrying. The style of "Prince" or "Princess" extends to the children of the Sovereign, the children of the sons of the Sovereign, and the eldest son of the eldest son of the Prince of Wales. Furthermore, wives of Princes are styled Princesses, though husbands of Princesses do not automatically become Princes.

UK Constitution and Government/Houses of Lancaster and York. Henry IV. Henry of Bolingbroke deposed his weak cousin, Richard II, in 1399. Henry IV's reign was marked by widespread rebellion. These were put down thanks to the great military skill of the Henry IV's son, the future King Henry V. Henry IV died in 1413 while plagued by a severe skin disease (possibly leprosy). Henry V. Henry V's reign was markedly different from his father's in that it involved little domestic turmoil. Overseas, Henry V's armies won several important victories in France. In 1415, the English defeated the French King Charles VI decisively at the Battle of Agincourt. About 100 English soldiers were killed, along with about 5000 Frenchmen. For the next two years, Henry V conducted delicate diplomacy to improve England's chances of conquering France. He negotiated with the Holy Roman Emperor Sigismund, who agreed to end the German alliance with France. In 1417, the war was renewed; by 1419, English troops were about to take Paris. The parties agreed to a treaty whereby Henry V was named heir of France. Henry V, however, died before he could succeed to the French throne, which therefore remained in the hands of the Frenchmen. Henry VI and Edward IV. Henry VI succeeded to the throne while still an infant. His uncles, the Dukes of Bedford and Gloucester, both functioned as Regents. During his reign, many French territories won during the Hundred Years War were lost. Henry VI's reign was interrupted by Edward IV's due to the War of the Roses. Henry VI was a member of the House of Lancaster, while Edward IV was from the House of York. The former House descended from Henry of Bolingbroke, the fourth son of King Edward III; the latter House descended from Edmund of Langley, Edward III's fifth son. In 1461, the Lancastrians lost to the Yorkists at the Battle of Towton. The Yorkist claimant, Edward IV, ascended to the throne, with the support of the powerful nobleman Richard Neville, 16th Earl of Warwick, known by the nickname "Warwick the Kingmaker". In 1464, Lancastrian revolts were put down. In 1469, however, Warwick the Kingmaker switched his allegiance, and in 1470, Henry VI was restored to the throne. The exiled Edward, however, soon returned and defeated Henry's forces. At the Battle of Tewkesbury, the remaining Lancastrians were defeated; Henry VI was also murdered. Edward V and Richard III. Edward IV was succeeded by his twelve year-old son in 1483. Edward IV's brother, Richard, was made guardian of Edward V and his brother, also named Richard. The young King's uncle usurped the throne and had Parliament declare the two brothers illegitimate. The two princes were then imprisoned in the Tower of London, where they might have been killed (their fate, however, is not certain). In 1485, Richard III faced Henry Tudor, the Lancastrian claimant, at the Battle of Bosworth Field, during which Richard became the last English monarch to be killed during battle. Henry came to power as Henry VII, establishing the Tudor Dynasty.

UK Constitution and Government/House of Tudor. Henry VII. Henry VII was one of the most successful monarchs in British history. He was the Lancastrian claimant to the throne and lived in France so as to remain safe from the designs of the Yorkist Kings. At the Battle of Bosworth Field in 1485, he defeated and killed the Yorkist Richard III. His claim was weak due to questions relating to the legitimacy of certain births, but he was nonetheless awarded the throne. Henry reformed the nation's taxation system and refilled the nation's treasury, which had been bankrupted by the fiscal irresponsibility of his predecessors. He also made peace with France so that the resources of the nation would not be spent trying to regain territories won during the Hundred Years' War. Henry also created marital alliances with Spain and Scotland. Henry's son, Arthur, married Catherine of Aragon, daughter of Ferdinand II of Aragon and Isabella I of Castile. Furthermore, Henry's daughter Margaret married James IV, King of Scots. When Henry's son Arthur died, he wished to protect the Anglo-Spanish alliance. Therefore, he obtained a dispensation from Pope Julius II allowing Henry's son, also named Henry, to marry Catherine. (Papal permission was necessary since Henry was marrying his brother's widow.) Upon Henry VII's death, Henry took the throne as Henry VIII. Henry VIII. King Henry VIII is often remembered for his multiple marriages. In his quest to obtain a male heir to the throne, Henry married six different times. His first marriage, as noted above, was to his brother's widow, Catherine of Aragon. That marriage occurred in 1509 and was scarred by several tragedies involving their children. The couple's first child was stillborn, their second lived for just 52 days, the third pregnancy ended as a miscarriage and the product of the fourth pregnancy died soon after birth. In 1516, the couple had a daughter, named Mary, followed by another miscarriage. Henry was growing impatient with his wife and eagerly sought a male heir. Henry sought to annul his marriage to Catherine. Ecclesiastic law permitted a man to marry his brother's widow only if the previous marriage had not been consummated. Catherine had informed the Pope that her marriage was non-consummate, so the Pope agreed to grant a dispensation allowing her to marry Henry. Now, however, Henry alleged that Catherine had lied, thereby rendering her marriage to him invalid. In 1533, an Act of Parliament annulled his marriage to Catherine, enabling him to marry Anne Boleyn. It was felt by many, however, that the Church, and not Parliament, could govern marriages. Henry had asked Pope Clement VII to issue a divorce several times. Under pressure from Catherine's nephew, Holy Roman Emperor Charles V, the Pope refused. Parliament therefore passed an Act denying appeals to Rome from certain decisions of English Archbishops. The Archbishop of Canterbury, Thomas Cranmer, annulled Henry's marriage to Catherine. In response, the Pope excommunicated Henry. Soon, the Church of England separated from the Roman Catholic Church. In 1534, all appeals to Rome from the decisions of the English clergy were stopped. An Act of Parliament passed in 1536 confirmed the King's position as "Supreme Head of the Church of England", thereby ending any ceremonial influence that the Pope still had. Anne Boleyn, meanwhile, was Henry's Queen, and the only surviving child from the marriage to Catherine, Mary, was declared illegitimate. Anne's first child, Elizabeth, was born in 1533. The next three pregnancies, however, all resulted in stillbirth or miscarriage. A dissatisfied Henry accused Anne of using witchcraft to entice him to marry her and to have five men enter into adulterous affairs with her. Furthermore, Anne was accused of treason because she had supposedly committed adultery while she was Queen. Anne's marriage to Henry VIII was annulled and she was executed at the Tower of London in 1536. Within two weeks of Anne's death, Henry married Jane Seymour. In 1537, Jane produced the male heir that Henry had long desired. The boy was named Edward and would later succeed Henry to the throne. Meanwhile, his half-sister Elizabeth was declared illegitimate. Shortly after the birth of the child, Jane died. Jane was followed as Queen by Anne of Cleeves, whom Henry married in 1540. Anne was the daughter of John III, Duke of Cleeves. Henry did not actually see Anne until shortly before their marriage; the relationship was contracted to establish an alliance between Henry and the Duke of Cleeves, a major Protestant leader. After Anne married him, Henry found her physically displeasing and unattractive. Shortly thereafter, the marriage was annulled on the grounds that Anne had previously been engaged to the Duke of Lorraine. After her divorce, Anne was treated well. She was given the title of Princess and allowed to live in Hever Castle, the former home of Anne Boleyn's family. Henry VIII's next marriage was to Catherine Howard, an Englishwoman of noble birth. In 1542, she was charged and convicted of high treason after having admitted to being engaged in an adulterous affair. In 1543, Henry contracted his final marriage, wedding Catherine Parr. The marriage lasted for the remainder of Henry's life, which ended in 1547. Edward VI and Lady Jane Grey. When Edward VI, son of Henry VIII and Jane Seymour, came to the throne, he was just ten years old. His uncle, Edward Seymour, Duke of Somerset served as Lord Protector while the King was a minor. Several nobles attempted to take over Somerset's role. John Dudley, 1st Earl of Warwick was successful; he was later created Duke of Northumberland. Edward VI was the first Protestant King of England. His father had broken away from the Roman Catholic Church but had not yet embraced Protestantism. Edward, however, was brought up Protestant. He sought to exclude his Catholic half-sister Mary from the line of succession. As he was dying at the age of fifteen, he made a document barring his half-sisters Mary and Elizabeth from the throne. He named the Lady Jane Grey, daughter-in-law of the Duke of Northumberland, his successor. Her claim to the throne was through her mother, who was a granddaughter of King Henry VII. Jane was proclaimed Queen upon Edward's death in 1553, but she served for only nine days before being deposed by Mary. Mary enjoyed far more popular support; the public also sympathised with the way her mother, Catherine of Aragon, had been treated. Jane was soon executed. She was seventeen years old at the time. Mary I. Mary was deeply opposed to her father's break from the Church in Rome. She sought to reverse reforms instituted by her Protestant half-brother. Mary even resorted to violence in her attempt to restore Catholicism, earning her the nickname "Bloody Mary". She executed several Protestants, including the former Archbishop of Canterbury Thomas Cranmer, on charges of heresy. In 1554, Mary married the Catholic King of Spain, Philip II. The marriage was unpopular in England, even with Catholic subjects. The couple were unable to produce a child before Mary's death from cancer in 1558. Elizabeth I. Mary's successor, her half-sister Elizabeth, was one of the most successful and popular British monarchs. The Elizabethan era was associated with cultural development and the expansion of English territory through colonialism. After coming to power, Elizabeth quickly reversed many of Mary's policies. Elizabeth reinstated the Church of England and had Parliament pass the Act of Supremacy, which confirmed the Sovereign's position as Supreme Governor of the Church of England. The Act also forced public and clerical officers to take the Oath of Supremacy recognising the Sovereign's position. Elizabeth, however, did practice limited toleration towards Catholics. After Pope Pius V excommunicated Elizabeth in 1570, Elizabeth ended her policy of religious toleration. One of Elizabeth's chief Catholic enemies was the Queen of Scotland, Mary. Since Elizabeth neither married nor bore any children, her cousin Mary was a possible heir to the English throne. Another possible heir was Lady Jane Grey's sister, Catherine. However, when Lady Catherine Grey died in 1568, Elizabeth was forced to consider that Catholic Mary was the most likely heir. Mary, however, had earlier been deposed by Scottish nobles, putting her infant son James on the throne. Mary had fled to England, hoping Elizabeth would aid her efforts to regain the Scottish throne, but Elizabeth reconsidered after learning of the "Ridolfi Plot", a scheme to assassinate Elizabeth and put the Roman Catholic Mary on the English throne. In 1572, Parliament passed a bill to exclude Mary from the line of succession, but Elizabeth refused to grant Royal Assent to it. Eventually, however, Mary proved to be too much of a liability due to her constant involvement in plots to murder Elizabeth. In 1587, she was executed after having been convicted of being involved in one such plot. Following Mary's execution, Philip II (widower of Mary I of England) sent a fleet of Spanish ships known as the "Armada" to invade England. England had supported a Protestant rebellion in the Netherlands and was seen as a threat to Catholicism. Furthermore, England had interfered with Spanish shipping and trade. Using Mary's execution as an excuse, Philip II obtained the Pope's authority to depose Elizabeth. In 1588, the Spanish Armada set sail for England. Harmed by bad weather, the Armada was defeated by Elizabeth's naval leaders, including Sir Francis Drake and the Lord Howard of Effingham. Towards the end of her life, Elizabeth still failed to name an heir. When she died, she was ironically succeeded by the son of Mary, Queen of Scots, James. James was already James VI, King of Scots; he became James I of England in 1603 and established the rule of the Stuart dynasty.

UK Constitution and Government/House of Stuart and the Commonwealth. James I. With the death of Elizabeth in 1603, the Crowns of England and Scotland united under James I. In 1567, when he was just a year old, James' mother Mary was forced to abdicate, and James became King James VI. Despite his mother's Catholicism, James was brought up as a Protestant. One of James' first acts as King was to conclude English involvement in the Eighty Years' War, also called the Dutch Revolt. Elizabeth had supported the Protestant Dutch rebels, providing one cause for Philip II's attack. In 1604, James signed the Treaty of London, thereby making peace with Spain. James had significant difficulty with the English Parliamentary structure. As King of Scots, he had not been accustomed to criticism from the Parliament. James firmly believed in the "Divine Right of Kings"—the right of Kings to rule that supposedly came from God—so he did not easily react to critics in Parliament. Under English law, however, it was impossible for the King to levy taxes without Parliament's consent, so he had to tolerate Parliament for some time. King James died in 1625 and was succeeded by his son Charles. Charles I. King Charles ruled at a time when Europe was moving toward domination by absolute monarchs. The French ruler, Louis XIV, epitomised this absolutism. Charles, sharing his father's belief in the Divine Right of Kings, also moved toward absolutist policies. Charles conflicted with Parliament over the issue of the Huguenots, French Protestants. Louis XIV had begun a persecution of the Huguenots; Charles sent an expedition to La Rochelle to provide aid to the Protestant residents. The effort, however, was disastrous, prompting Parliament to further criticise him. In 1628, the House of Commons issued the Petition of Right, which demanded that Charles cease his use of arbitrary power. Charles had persecuted individuals using the Court of the Star Chamber, a secret court that could impose any penalty, even torture, except for death. Charles had also imprisoned individuals without a trial and denied them the right to the writ of "habeas corpus". The Petition of Right, however, was not successful; in 1629, Charles dissolved Parliament. He ruled alone for the next eleven years, which is sometimes referred to as the "eleven years of tyranny" or "personal rule". Since Parliamentary approval was required to impose taxes, Charles had grave difficulty in keeping the government functional. Charles imposed several taxes himself; these were widely seen as unlawful. During these eleven years, Charles began instituting religious reforms in Scotland, moving it towards the English model. He attempted to impose the Anglican Prayer Book on Scottish churches, leading to riots and violence. In 1638, the General Assembly of the Church of Scotland abolished the office of bishop and established Presbyterianism (an ecclesiastic system without clerical officers such as bishops and archbishops). Charles sent his armies to Scotland, but was quickly forced to end the conflict, known as the First Bishops' War, because of a lack of funding. Charles granted Scotland certain parliamentary and ecclesiastic freedoms in 1639. In 1640, Charles finally called a Parliament to authorise additional taxation. Since the Parliament was dissolved within weeks of its summoning, it was known as the "Short Parliament". Charles then sent a new military expedition to Scotland to fight the Second Bishops' War. Again, the Royal forces were defeated. Charles then summoned Parliament again, this Parliament becoming known as the "Long Parliament", in order to raise funds for making reparations to the Scots. Tension between Charles and Parliament increased dramatically. Charles agreed to abolish the hated Star Chamber, but he refused to give up control of the army. In 1641, Charles entered the House of Commons with armed guards in order to arrest his Parliamentary enemies. They had already fled, however, and Parliament took the breach of their premises very seriously. (Since Charles, no English monarch has sought to set foot in the House of Commons.) The unsafe monarch moved the Royal court to Oxford. Royal forces controlled north and west England, while Parliament controlled south and east England. A Civil War broke out, but was indecisive until 1644, when Parliamentary forces clearly gained the upper hand. In 1646, Charles was forced to escape to Scotland, but the Scottish army delivered him to Parliament in 1647. Charles was then imprisoned. Charles negotiated with the Scottish army, declaring that if it restored him to power, he would implement the Scottish Presbyterian ecclesiastic model in England. In 1648, the Scots invaded England, but were defeated. The House of Commons began to pass laws without the consent of either the Sovereign or the House of Lords, but many MPs still wished to come to terms with the king. Members of the army, however, felt that Charles had gone too far by siding with the Scots against England and were determined to have him brought to trial. In December 1648 an army regiment, Colonel Pride's, used force to bar entry into the House of Commons, only allowing MPs who would support the army to remain. These MPs, the "Rump Parliament", established a commission of 135 to try Charles for treason. Charles, an ardent believer in the Divine Right of Kings, refused to accept the jurisdiction of any court over him. Therefore, he was by default considered guilty of high treason and was executed on January 30, 1649. Oliver and Richard Cromwell. At first, Oliver Cromwell ruled along with the republican Parliament, the state being known as the "Commonwealth of England". After Charles' execution, however, Parliament became disunited. In 1653, he suspended Parliament, and as Charles had done earlier, began several years of rule as a dictator. Later, Parliament was recalled, and in 1657 offered to make Cromwell the King. Since he faced opposition from his own senior military officers, Cromwell declined. Instead, he was made a "Lord Protector", even being installed on the former King's throne. He was a King in all but name. Cromwell died in 1658 and was succeeded by his son Richard, an extremely poor politician. Richard Cromwell was not interested in his position and abdicated quickly. The Protectorate was ended and the Commonwealth restored. Anarchy was the result. Quickly, Parliament chose to reestablish the monarchy by inviting Charles I's son to take the throne as Charles II. Charles II. During the rule of Oliver Cromwell, Charles II remained King in Scotland. After an unsuccessful challenge to Cromwell's rule, Charles escaped to Europe. In 1660, when England was in anarchy, Charles issued the Declaration of Breda, outlining his conditions for returning to the Throne. The Long Parliament, which had been convened in 1640, finally dissolved itself. A new Parliament, called the "Convention Parliament", was elected; it was far more favourable to the Royalty than the Long Parliament. In May 1660, the Convention Parliament that Charles had been the lawful King of England since the death of his father in 1649. Charles soon arrived in London and was restored to actual power. Charles granted a general pardon to most of Cromwell's supporters. Those who had directly participated in his father's execution, however, were either executed or imprisoned for life. Cromwell himself suffered a posthumous execution: his body was exhumed, hung, drawn and quartered, his head cut off and displayed from a pole and the remainder of his body thrown into a common pit. The posthumous execution took place on the anniversary of Charles I's death. Charles also dissolved the Convention Parliament. The next Parliament, called the "Cavalier Parliament" was soon elected. The Cavalier Parliament lasted for seventeen years without an election before being dissolved. During its long tenure, the Cavalier Parliament enacted several important laws, including many that suppressed religious dissent. The Act of Uniformity required the use of the Church of England's "Book of Common Prayer" in all Church services. The Conventicle Act prohibited religious assemblies of more than five members except under the Church of England. The Five Mile Act banned non-members of the Church of England from living in towns with a Royal Charter, instead forcing them into the country. In 1672, Charles mitigated these laws with the Royal Declaration of Indulgence, which provided for religious toleration. Parliament, however, suspected him of Catholicism and forced him to withdraw the Declaration. In 1673, Parliament passed the Test Act, which required civil servants to swear an oath against Catholicism. Parliament's suspicions did turn out to be accurate. As Charles II lay dying in 1685, he converted to Catholicism. Charles did not have a single legitimate child, though he did have, while living in Europe, several illegitimate ones (over 300 by some estimates). He was succeeded, therefore, by his younger brother James, an open Catholic. James II. James II (James VII in Scotland) was an extremely controversial monarch due to his Catholicism. Soon after he took power, a Protestant illegitimate son of Charles II, James Scott, Duke of Monmouth, proclaimed himself King. James II defeated him within a few days and had him executed. James made himself highly unpopular by appointing Catholic officials, especially in Ireland. Later, he established a standing army in peacetime, alarming many Protestants. Rebellion, however, did not occur because people trusted James' daughter Mary, a Protestant. In 1688, however, James produced a son, who was brought up Catholic. Since Mary's place in the line of succession was lowered, and a Catholic Dynasty in England seemed ineveitable, the "Immortal Seven"—the Duke of Devonshire, the Earl of Danby, the Earl of Shrewsbury, the Viscount Lumley, the Bishop of London, Edward Russell and Henry Sidney—conspired to replace James and his son with Mary and her Dutch husband William of Orange. In 1688, William and Mary invaded England and James fled the country. The revolution was hailed as the "Glorious Revolution" or the "Bloodless Revolution". Though the latter term was inaccurate, the revolution was not as violent as the War of the Roses or the English Civil War. William and Mary. Parliament wished then to make Mary the sole Queen. She, however, refused and demanded that she be made co-Sovereign with her husband. In 1689, the Parliament of England declared in the English Bill of Rights, one of the most significant constitutional documents in British history, that James' flight constituted an abdication of the throne and that the throne should go jointly to William (William III) and Mary (Mary II). The Bill of Rights also required that the Sovereign cannot deny certain rights, such as freedom of speech in Parliament, freedom from taxation without Parliament's consent and freedom from cruel and unusual punishment. In Scotland, the Estates General passed a similar Act, called the Claim of Right, which also made William and Mary joint rulers. In Ireland, power had to be won in battle. In 1690, the English won the Battle of the Boyne, thereby establishing William and Mary's rule over the entire British Isles. For the early part of the reign, Mary administered the Government while William controlled the military. Unpopularly, William appointed people from his native Holland as officers in the English army and Royal Navy. Furthermore, he used English military resources to protect the Netherlands. In 1694, after the death of Queen Mary from smallpox, William continued to rule as the sole Sovereign. Since William and Mary did not have children, William's heir was Anne, who had seventeen pregnancies, most of which ended in stillbirth. In 1700, Anne's last surviving child, William, died at the age of eleven. Parliament was faced with a succession crisis, because after Anne, many in the line of succession were Catholic. Therefore, in 1701, the Act of Settlement was passed, allowing Sophia, Electress and Duchess Dowager of Hanover (a German state), and her Protestant heirs, to succeed if Anne had no further children. Sophia's claim stemmed from her great-grandfather, James I. Several lines that were more senior to Sophia's were bypassed under the act. Some of these had questionable legitimacy, while others were Catholic. The Act of Settlement also banned non-Protestants and those who married Catholics from the throne. In 1702, William died, and his sister-in-law Anne became Queen. Anne. Even following the passage of the Act of Settlement, Protestant succession to the throne was insecure in Scotland. In 1703, the Scottish Parliament, the Estates, passed a bill that required that, if Anne died without children, the Estates could appoint any Protestant descendant of Scottish monarchs as the King. The individual appointed could not be the same person who would, under the Act of Settlement, succeed to the English crown unless several economic conditions were met. The Queen's Commissioner refused Royal Assent on her behalf. The Scottish Estates then threatened to withdraw Scottish troops from the Queen's armies, which were then engaged in the War of the Spanish Succession in Europe and Queen Anne's War in North America. The Estates also threatened to refuse to levy taxes, so Anne relented and agreed to grant Royal Assent to the bill, which became the Act of Security. The English Parliament feared the separation of the Crowns which had been united since the death of Elizabeth I. They therefore attempted to coerce Scotland, passing the Alien Act in 1705. The Alien Act provided for cutting off trade between England and Scotland. Scotland was already suffering from the failure of the Darién Scheme, a disastrous and expensive attempt to establish Scottish colonies in America. Scotland quickly began to negotiate union with England. In 1707, the Act of Union was passed, despite mass protest in Scotland, by Parliament and the Scottish Estates. The Act combined England and Scotland into one Kingdom of Great Britain, terminated the Parliament and Estates, and replaced them with one Parliament of Great Britain. Scotland was entitled to elect a certain number of members of the House of Commons. Furthermore, it was permitted to send sixteen of its peers to sit along with all English peers in the House of Lords. The Act guaranteed Scotland the right to retain its distinct legal system. The Church of Scotland was also guaranteed independence from political interference. Ireland remained a separate country, though still governed by the British Sovereign. Anne is often remembered as the last British monarch to deny Royal Assent to a bill, which she did in 1707 to a militia bill. Due to her poor health, made worse by her failed pregnancies, her government was run through her ministers. She died in 1714, to be succeeded by George, Elector of Hanover, whose mother Sophia had died a few weeks earlier.

UK Constitution and Government/House of Hanover. George I. George, Duke and Elector of Hanover became King George I in 1714. His claim was opposed by the Jacobites, supporters of the deposed King James II. Since James II had died, his claim was taken over by his son, James Francis Edward Stewart, the "Old Pretender." In 1715, there was a Jacobite rebellion, but an ill James could not lead it. By the time he recovered, it was too late, and the rebellion was suppressed. King George was not deeply involved in British politics; instead, he concentrated on matters in his home, Germany. The King could not even speak English, earning the ridicule of many of his subjects. George, furthermore, spent much time in his native land of Hanover. Meanwhile, a ministerial system developed in Great Britain. George appointed Sir Robert Walpole as "First Lord of the Treasury". Walpole was George's most powerful minister, but he was not termed "Prime Minister"; that term came into use in later years. Walpole's tenure began in 1721; other ministers held office at his, rather than the King's, pleasure. George's lack of involvement in politics contributed greatly to the development of the modern British political system. George died in 1727 from a stroke while in Germany. He was succeeded by his son, who ruled as George II. George II. George II was naturalised as a British citizen in 1705; his reign began in 1727. Like his father, George transferred political power to Sir Robert Walpole, who served until 1742. Walpole was succeeded by Spencer Compton, 1st Earl of Wilmington, who served until 1743, and then by Henry Pelham, who served until his death in 1754. During Pelham's service, the nation experienced a second Jacobite Rebellion, which was almost successful in putting Bonnie Prince Charlie—son of the Old Pretender, himself called the Young Pretender—on the throne. The rebellion began in 1745 and was ended in 1746 when the King's forces defeated the Jacobites at the Battle of Culloden, the last battle ever to be fought on British soil. Before George II's death in 1760, he was served by two other Prime Ministers: Henry Pelham's elder brother the Duke of Newcastle, and the Duke of Devonshire. George II's eldest son, Frederick, had predeceased him, so George was succeeded by his grandson, also named George. George III. George III attempted to reverse the trend that his Hanoverian predecessors had set by reducing the influence of the Prime Minister. He appointed a variety of different people as his Prime Minister, on the basis of favouritism rather than ability. The Whig Party of Robert Walpole declared George an autocrat and compared him to Charles I. George III's reign is notable for many important international events. In 1763, Great Britain defeated France in the Seven Years' War, a global war that also involved Spain, Portugal and the Netherlands and was fought in Europe, America and India. As a result of the Treaty of Paris, New France (the French territory in North America, including Quebec and land east of the Mississippi) was ceded to Britain, as was Spanish Florida. Spain, however, took New Orleans and Louisiana, the vast French territory on the west of the Mississippi. Great Britain came to be recognised as the world's pre-eminent colonial power, displacing France. The nation, however, was left deeply in debt. To overcome it, British colonies in America were taxed, much to their distaste. Eventually, Britain lost its American colonies during the American War of Independence, which lasted from 1776 to 1783. Elsewhere, however, the British Empire continued to expand. In India, the British East India Company took control of many small nation-states nominally headed by their own princes. The island of Australia was also occupied, and Canada's population increased with the number of British Loyalists who left the newly formed United States of America. In 1801, Parliament passed the Act of Union, uniting Great Britain and Ireland into the United Kingdom. Ireland was allowed to elect 100 Members of Parliament to the House of Commons and 22 representative peers to the House of Lords. The Act originally provided for the removal of restrictions from Roman Catholics, but George III refused to agree to the proposal, arguing that doing so would violate his oath to maintain Protestantism. George was the last British monarch to claim the Kingdom of France. He was persuaded to abandon the meaningless claim dating to the Plantagenet days in 1801 by the French ruler Napoleon. In 1811, George III, who had previously suffered bouts of madness, went permanantly insane. His son George ruled the country as "Prince Regent", and became George IV when the King died in 1820. George IV. George IV is often remembered as an unwise and extravagant monarch. During his Regency, London was redesigned, and funding for the arts was increased. As King, George was unable to govern effectively; he was overweight, possibly addicted to a form of opium and showing signs of his father's mental disease. While he ruled, George's ministers were once again able to regain the power that they had lost during his father's reign. George opposed several popular social reforms. As his father, he refused to lift several restrictions on Roman Catholics. Upon his death in 1830, his younger brother began to reign as William IV. William IV. Early in William's reign, British politics was reformed by the Reform Act of 1832. At the time, the House of Commons was a disorganised and undemocratic body, unlike the modern House. The nation included several "rotten boroughs", which historically had the right to elect members of Parliament, but actually had very few residents. The rotten borough of Old Sarum, for instance, had seven voters, but could elect two MPs. An even more extreme example is of Dunwich, which could also elect two MPs despite having no residents, the entire borough having been eroded away into the North Sea. Other boroughs were called "pocket boroughs" because they were "in the pocket" of a wealthy landowner, whose son was normally elected to the seat. At the same time, entire cities such as Westminster (with about 20,000 voters) still had just two MPs. The House of Commons agreed to the Reform Bill, but it was rejected by the House of Lords, whose members controlled several pocket boroughs. The Tory Party, furthermore, opposed the bill actively. William IV agreed with his Prime Minister, the Earl Grey, to flood the House of Lords with pro-reform members by creating fifty new peerages; when the time came, he backed down. The Earl Grey and his Whig Party government then resigned, but returned to power when William finally agreed to co-operate. The Reform Act of 1832 gave urban areas increased political power, but allowed aristocrats to retain effective control of the rural areas. Over fifty rotten boroughs were abolished, while the representation of some other boroughs was reduced from two MPs to one. Though members of the middle class were granted the right to vote, the Reform Act did not do much to expand the electorate, which amounted after passage to just three percent of the population. In 1834, William became the last British monarch to appoint a Prime Minister who did not have the confidence of Parliament. He replaced the Whig Prime Minister, the Viscount Melbourne, with a Tory, Sir Robert Peel. Peel, however, had a minority in the House of Commons, so he resigned in 1835, and Melbourne returned to power. In 1837, William died and was succeeded on the British throne by his niece Victoria, who was just eighteen years old at the time. The union of the Crowns of Britain and Hanover was then dissolved, since Salic Law, which applied in Hanover, only allowed males to rule. Therefore, Hanover passed to William's brother Ernest. Victoria. A few years after taking power, Victoria married a German Prince, Albert of Saxe-Coburg-Gotha, who was given the title of "Prince Consort". Albert originally wished to actively govern the United Kingdom, but he acquiesced to his wife's requests to the contrary. The extremely happy marriage ended with Albert's death in 1861, following which Victoria entered a period of semi-mourning that would last for the rest of her reign. She was often called "the Widow of Windsor", after Windsor Castle, a Royal home. In 1867, Parliament passed another Reform Act. Like its predecessor, the Reform Act of 1832, true electoral reform was not achieved; the property qualifications limited the electorate to about eight percent of the population. Therafter, power was held by two Prime Ministers—Benjamin Disraeli (a Tory and a favourite of Victoria) and William Ewart Gladstone (a Liberal whom Victoria disliked)—from 1868 to 1885. In 1876, Disraeli convinced Victoria to take the title of Empress of India. Many of Victoria's daughters married into European Royal Houses, giving her the nickname "Grandmother of Europe". All of the current European monarchs descend from Victoria. Victoria died in 1901, holding the record for longest serving British Sovereign. She was succeeded by her son Edward, who became King Edward VII. Edward was deemed to belong not to his mother's House of Hanover, but instead to his father's dynasty, Saxe-Coburg-Gotha.

UK Constitution and Government/Houses of Saxe-Coburg-Gotha and Windsor. Edward VII. Edward VII was the oldest person in British history to become King, beginning his reign at the age of fifty-nine. He participated actively in foreign affairs, visiting France in 1903. The visit led to the "Entente Cordiale" (Friendly Understanding), an informal agreement between France and the United Kingdom marking the end of centuries of Anglo-French rivalry. In the case of Germany, however, Edward VII exacerbated rivalry through his bad relations with his nephew, Kaiser Wilhelm II. Towards the end of his life, Edward was faced with a constitutional crisis when the Liberal Government, led by Herbert Henry Asquith, proposed the "People's Budget". The Budget reformed the tax system by creating a land tax, which would adversely affect the aristocratic class. The Conservative landowning majority in the House of Lords broke convention by rejecting the budget. They argued that the Commons themselves had broken a convention by attacking the wealth of the Lords. Before the problem could be resolved, Edward VII died in 1910, allowing his son, George, to ascend to the throne. George V. After George became King, the constitutional crisis was resolved after the Liberal Government resigned and Parliament was dissolved. The Liberals were reelected, in part due to the unpopularity of the House of Lords, and used the election as a mandate to force their Budget through, almost too late to save the nation's financial system from ruin. The Lords paid a price for their opposition to the Liberals, who in the commons passed the Parliament Bill, which provided that a bill could be submitted for the King's Assent if the Commons passed it in three consecutive sessions, even if the Lords rejected it. The time would later be reduced to two sessions in 1949. When the House of Lords refused to pass the Parliament Bill, Prime Minister Asquith asked George V to create 250 new Liberal peers to erase the Conservative majority. George agreed, but the Lords acquiesced and passed the bill quickly. World War I occurred during George's reign. Due to the family's German connections, the Royalty began to become unpopular; George's cousin, Wilhelm II, was especially despised. In 1917, to appease the public, George changed the Royal House's name from the German-sounding "Saxe-Coburg-Gotha" to the more English "Windsor". In 1922, most of Ireland left the United Kingdom to form the Irish Free State following the Irish Civil War. The Irish Free State retained the British monarch as a Sovereign, but functioned as a Dominion of the Crown, with its own Government and Legislature. Six counties in the Irish province of Ulster remained in the United Kingdom as Northern Ireland. In 1927, the name of the country was changed from "the United Kingdom of Great Britain and Ireland" to "the United Kingdom of Great Britain and Northern Ireland". George V died in 1936 and was succeeded by his son, who ruled as Edward VIII. Edward VIII. Edward VIII became King in January of 1936 and abdicated in December. His reign was controversial because of his desire to marry the American Wallis Simpson. Simpson was already divorced once; she divorced her second husband so she could marry King Edward. A problem, however, existed because Edward was the Supreme Governor of the Church of England, which prohibited remarriage after divorce. The Government advised him that he could not marry while he was King, so he indicated a desire to abdicate and marry Simpson. The abdication was not unilateral, as the Act of Settlement provided that the Crown go to the heir of Sophia, Electress of Hanover, regardless of that person's willingness to rule. Therefore, Parliament had to pass a special Act in order to permit Edward to abdicate, which he did. Edward's brother, Albert Frederick Arthur George, became King. He chose to rule as George VI to create a link in the public's mind between him and the previous Kings of the same name during a time of crisis. Edward, meanwhile, was made Duke of Windsor and the issue of his marriage to Simpson were excluded from the line of succession. George VI. When George took power in 1936, the popularity of the Royal Family had been damaged by the abdication crisis. It was, however, restored when George and his wife, Queen Elizabeth, led the nation and boosted morale during World War II. During the war, Britain was led by one of its most famous Prime Ministers, Sir Winston Churchill. Following the War, the United Kingdom began to lose several of its overseas possessions. In 1947, India became independent and George lost the title of Emperor of India. Until 1950, however, he remained King of India while a constitution was being written. George was also the last King of Ireland; the Irish established a republic in 1949. George died in 1952 from lung cancer. His daughter Elizabeth succeeded him. Elizabeth II. During Elizabeth's reign, there have been several important constitutional developments. A notable one occurred in 1963, when Conservative Prime Minister Harold Macmillan resigned. There was no clear leader of the Conservative Party, but many favoured Richard Austen Butler, the Deputy Prime Minister. Harold Macmillan advised the Queen, however, that senior politicians in the party preferred Alec Douglas-Home, 14th Earl of Home. Elizabeth accepted the advice and appointed the Earl of Home to the office of Prime Minister, marking the last time a member of the House of Lords would be so appointed. Home, taking advantage of the Peerage Act passed in 1963, "disclaimed" his peerage. A Conservative member of the House of Commons vacated his seat, allowing Home to contest the by-election for that constituency and become a member of the House of Commons. There have also been many recent constitutional developments in the nation. The office of Prime Minister increased greatly in power under the Conservative Prime Minister Margaret Thatcher (the "Iron Lady") and the Labour Prime Minister Tony Blair. Under Blair, many of Parliament's lawmaking functions were devolved to local administrations in Scotland, Wales and Northern Ireland. In 1999, the House of Lords Act was passed, removing the automatic right of hereditary peers to sit in the House. Charles III. Charles III acceeded to the throne, aged 74, on the death of his mother, Elizabeth II in 2022 and was crowned in 2023.

UK Constitution and Government/Government. Structure. "His Majesty's Government" is the executive political authority for the United Kingdom as a whole. At its heart is the "Cabinet", a grouping of senior "Ministers of the Crown", headed by the "Prime Minister". Members of the Government are political appointees, and are usually drawn from one of the two Houses of Parliament. In addition to the heads of the "Departments of State" (most of whom carry the title of "Secretary of State"), the Government also includes "junior ministers" (who bear the title of "Under Secretary of State", "Minister of State", or "Parliamentary Secretary"), "whips" (responsible for enforcing party discipline within the two Houses), and "Parliamentary Private Secretaries" (political assistants to ministers). When the Sovereign is the Queen, the Government is referred to as "Her" Majesty's Government; likewise, when there are joint Sovereigns, the Government is known as "Their Majesties"' Government. Prime Minister. The Prime Minister (or "PM") is the head of the Government. Since the early twentieth century the Prime Minister has held the office of "First Lord of the Treasury", and in recent decades has also held the office of "Minister for the Civil Service". The Prime Minister is "asked to form a Government" by the Sovereign. Usually this occurs after a general election has altered the balance of party political power within the House of Commons. The Prime Minister is expected to "have the confidence" of the House of Commons; this usually means that he or she is the leader of the political party holding the majority of the seats in the Commons. Since at least the 1920s the Prime Minister himself is also expected to be a Member of Parliament (i.e. member of the House of Commons). The Prime Minister retains office until he or she dies or resigns, or until someone else is appointed; this means that even when expecting to be defeated at a general election, the Prime Minister remains formally in power until his or her rival is returned as an MP and asked in turn to form a Government. By law, if defeated by an actual Commons vote of "no confidence", there is a period of 14 days during which an alternative Government may be formed, if a "vote of confidence" can be won by a prospective Prime Minister. If no Government can be formed within the 14 days an election is automatically triggered, in effect allowing the electorate itself to approve or disapprove of the Prime Minister's policy. The polling day for the election is to be the day appointed by Her Majesty by proclamation on the recommendation of the Prime Minister. The Parliament then in existence dissolves at the beginning of the 17th working day before the polling day. The existence and basis of appointment of the office is a matter of constitutional convention rather than of law. Because of this, there are no formal qualifications for the office. However, a small number of Acts of Parliament do make reference to the Prime Minister, and since the 1930s office has carried a salary in its own right. The Prime Minister is often an extremely powerful figure within the political system; the office has been said by some to be an "elected dictatorship", and some Prime Ministers have been accused of being "presidential". A weak Prime Minister may be forced out of office (i.e. forced to resign) by his or her own party, particularly if there is an alternative figure within the party seen as a better choice. Cabinet and other ministers. Membership of the Cabinet is not defined by law, and is only loosely bound by convention. The Prime Minister and (if there is one) the Deputy Prime Minister are always members, as are the three most senior ministerial heads of Departments of State: the "Secretary of State for Foreign and Commonwealth Affairs" (commonly known as the "Foreign Secretary"), the "Chancellor of the Exchequer" (i.e. the minister responsible for finance), and the "Secretary of State for the Home Department" (commonly known as the "Home Secretary"). Most of the other heads of departments are usually members of the Cabinet, as well as a small number of junior ministers. Ministers of the Crown are formally appointed by the Sovereign upon the "advice" of the Prime Minister. Ministers are bound by the convention of "collective responsibility", by which they are expected to publicly support or defend the policy of the Government, or else resign. They are also bound by the less clearly defined convention of "individual responsibility", by which they are responsible to Parliament for the acts of their department. Ministers are often called upon to resign who either by their own actions, or by those of their department, are perceived in some manner to have failed in their duty; however, it usually takes sustained criticism over a period of time for both a minister to feel compelled to resign, and for the Prime Minister to accept that resignation. Occasionally a minister offers his or her resignation, but the Prime Minister retains them in office. Parliamentary Private Secretaries are also bound by the principle of "collective responsibility", even though they hold no ministerial responsibility and take no part in the formation of policy; the position is seen as an initial stepping-stone towards being offered ministerial office. Privy Council. "His Majesty's Most Honourable Privy Council" is a ceremonial body of advisors to the Sovereign. The Privy Council is used as a mechanism for maintaining ministerial responsibility for the actions of the Crown; for example, royal proclamations are approved by the Privy Council before they are issued. All senior members of the Government are appointed to be Privy Counsellors, as well as certain senior members of the Royal Family, leaders of the main political parties, the archbishops and senior bishops of the Church of England, and certain senior judges. The Privy Council is headed by the "Lord President of the Council", a ministerial office usually held by a member of the Cabinet. By convention the Lord President is also either the "Leader of the House of Commons", or the "Leader of the House of Lords", with responsibility for directing and negotiating the course of business in the respective House. Meetings of the Privy Council are usually extremely short, and are rarely attended by more than a bare minimum of Privy Counsellors.

UK Constitution and Government/Judiciary. Structure. The United Kingdom is made up of three separate legal jurisdictions, each with a separate laws and hierarchy of courts: "England and Wales", "Scotland", and "Northern Ireland". England and Wales. England and Wales is a "common law" jurisdiction. Lower courts. The lowest court in England and Wales is the "Magistrates' Court". Magistrates, also known as Justices of the Peace, are laypersons appointed by the Sovereign. The court hears "summary" offences (punishable by six months or less in prison). When hearing such cases, three magistrates sit together as a panel without a jury. In some metropolitan areas, such as London, there are no magistrates; instead, summary cases are tried by a single District Judge who is trained in law. Serious criminal cases are tried before a "Crown Court" with a judge and a jury of twelve. The accused may also choose to have certain summary offences referred from the magistrates' court to the Crown Court, in order for their case to be tried before a jury; the Crown Court also hears appeals from magistrates' courts. Though the Crown Court is constituted as a single body for the whole of England and Wales, it sits permanently at multiple places throughout its area of jurisdiction. The counterpart to the Crown and Magistrates' Courts in the civil justice system is the "County Court". There are over 200 County Courts throughout England and Wales. High Court. The "High Court of England and Wales" takes appeals from the County Court, and also has an original jurisdiction in certain matters. The High Court is constituted into three "divisions": the "Family Division", the "Chancery Division", and the "Queen's Bench Division". The Family Division is presided over by the "President of the Family Division", and hears cases involving family matters such as matrimonial breakdown, child custody and welfare, and adoption. The Chancery Division is presided over by the "Chancellor of the High Court" (formerly known as the "Vice-Chancellor"), and hears cases involving land, companies, bankruptcy, and probate. The Queen's Bench Division is presided over by the "President of the Queen's Bench Division", and hears cases involving torts (civil wrongs). The Queen's Bench Division also includes four subordinate courts: the "Admiralty Court" (dealing with shipping), the "Commercial Court" (dealing with insurance, banking, and commerce), the "Technology and Construction Court" (dealing with complex technological matters), and the "Administrative Court" (exercising judicial review over the actions of local government). The Queen's Bench Division also has oversight of the lower courts. Court of Appeal. Above the High Court in civil cases, and the Crown Court in criminal cases, is the Court of Appeal, headed by the "Master of the Rolls", and including 35 "Lords Justices of Appeal" as well as other judges. The Court of Appeal is divided into a "Civil Division" (presided over by the Master of the Rolls) and a "Criminal Division" (presided over by the Lord Chief Justice). Generally speaking, appeals may only be heard "by leave"; that is, with the permission of the either the Court of Appeal or the judge whose decision is being contested. In some cases, it is possible to "leapfrog" the High Court and bring a case directly from a County Court. Together, the Crown Court, the High Court, and Court of Appeal constitute the "Senior Courts" (formerly known as the "Supreme Court of Judicature"). Thus, since they are theoretically one body, it is possible for judges of one court to sit in other courts. Appeals from the Senior Courts go to the Supreme Court; it is also possible to leapfrog from the High Court, but not from the Crown Court. Normally, leave to appeal to the Supreme Court is not granted unless the case is of great legal or constitutional importance. Northern Ireland. Northern Ireland's system is based on that used in England and Wales, with a similar hierarchy of magistrates' court, the Crown Court (for criminal trials), county courts (for civil trials), the High Court, and the Court of Appeal. Appeals from Northern Ireland lie to the Supreme Court. Scotland. In contrast with the rest of the United Kingdom, Scotland uses a mixture of common law and civil law. Its court system was developed independently of that in England. The "Act of Union" (1707) guarantees the continuance of Scotland's different legal system. Lower courts. Summary jurisdiction is exercised by "Justice of the Peace Courts", held either by three "Justices of the Peace" (lay magistrates) sitting together, or by a Justice of the Peace sitting with a legally qualified clerk. As in England and Wales, professional judges may sit in certain metropolitan areas. Above the Justice of the Peace Courts are the "Sheriff Courts", of which there are around 50. Sheriff Courts hear both criminal and civil cases, and are held before a judge known as a Sheriff, and have a jury of fifteen people. Sheriff Courts are grouped into six different Sheriffdoms, headed by a Sheriff Principal who hears appeals from cases not decided by a jury. High Court of Justiciary. The highest criminal court in Scotland is the "High Court of Justiciary". The judges of the court are also the judges of the Court of Session (see below); as High Court judges they are known as "Lords Commissioners of Justiciary". The head of the court is the "Lord Justice-General" (also the Lord President of the Court of Session), with a deputy known as the "Lord Justice Clerk" (who holds the same office in the Court of Session). Altogether the High Court has up to 32 individual judges. The High Court has exclusive jurisdiction in serious crimes, such as murder or drug trafficking, in which case a single judge sits with a jury of fifteen. The High Court also hears appeals from Justice of the Peace Courts, and hears appeals in criminal cases from Sheriff Courts. Appeals against decisions by a High Court judge in criminal cases are heard by either two (in appeals against sentences) or three (in appeals against conviction) High Court judges. No appeal lies beyond the High Court. Court of Session. The highest civil court in Scotland is the "Court of Session". Its judges also sit as judges of the High Court of Justiciary (see above); as Court of Session judges they are known as "Lords and Ladies of Council and Session", or "Senators of the College of Justice". The Court is headed by the "Lord President", with a "Lord Justice Clerk" as deputy. Altogether the Court of Session has up to 32 individual judges. The Court of Session is divided into the "Outer House" (made up of nineteen judges), and the "Inner House" (made up of the remaining judges). The Outer House has original jurisdiction, while the Inner House has appellate jurisdiction. The Inner House is further divided into the First and Second Divisions, headed by the Lord President and Lord Justice Clerk respectively. Sometimes, when many cases are before the court, an Extra Division may be appointed. Each Division may sit as a panel hearing an appeal from the Sheriff Court or from the Outer House. Appeals from the Court of Session lie to the Supreme Court. Supreme Court. The "Supreme Court of the United Kingdom" is the ultimate court of appeal in all civil matters, as well as in criminal cases (other than from Scotland), and also has original jurisdiction in devolution cases. The Supreme Court has replaced the jurisdiction previously exercised by the House of Lords in the latter's now-abolished judicial capacity. The Supreme Court of the United Kingdom is not to be confused with the Supreme Court of Judicature, the name formerly held by (a) the Senior Courts, in England and Wales, and (b) the Court of Judicature, in Northern Ireland. The Supreme Court is headed by a President, who has a Deputy President. There are a further ten "puisne" judges. Judicial Committee of the Privy Council. The "Judicial Committee of the Privy Council" formerly held original jurisdiction in the United Kingdom in devolution cases, and continues to hold appellate jurisdiction over the ecclesiastical courts of the Church of England. Appeals to the Privy Council as a court of last resort also lie from the Crown dependencies, the British overseas territories, and from certain Commonwealth countries. Membership of the Judicial Committee is made up of Justices of the Supreme Court, Privy Counsellors who are or were Lord Justices of Appeal in either England and Wales or Northern Ireland, members of the Inner House of Scotland's Court of Session, and selected senior judges from certain other Commonwealth countries. Members retire at the age of 75. Appeals to "Her Majesty in Council" are referred to the Judicial Committee, which formally reports to the Queen in Council, who in turn formally confirms the report. By agreement, appeals from certain Commonwealth countries lie directly to the Judicial Committee itself. The Queen-in-Council also considers appeals from the disciplinary committees of certain medical bodies such as the Royal College of Surgeons. Also, cases against the Church Commissioners (who administer the Church of England's property estates) may be considered. Appeals may be heard from certain ecclesiastic courts (the Court of Arches in Canterbury, and the Chancery Court in York) in cases that do not involve Church doctrine. Appeals may also be heard from certain dormant courts, including Prize Courts (which hear cases relating to the capture of enemy ships at sea, and the ownership of property seized from captured ships) and the Court of Admiralty of the Cinque Ports. Finally, the Queen-in-Council determines if an individual is qualified to be elected to the House of Commons under the House of Commons Disqualification Act. ECHR and ECJ. In addition to the above domestic courts, there are two further courts which can be said to exercise a jurisdiction over the United Kingdom. The "European Court of Human Rights" deals with cases concerning alleged infringements of the "European Convention on Human Rights". The "European Court of Justice" deals with cases concerning alleged infringements of European Union law.

UK Constitution and Government/Constitution. Unlike most other sovereign states, the United Kingdom does not possess a document expressing itself to be the nation's fundamental or highest law. Instead, the British constitution is found in a number of sources. Because of this, the British constitution is often said to be an "unwritten constitution"; however, many parts of the constitution are indeed in written form, so it would be more accurate to refer to the body of the British constitution as an "uncodified constitution". The British constitution is spread across a number of sources: 1. Statute law 2. Royal prerogative (executive powers usually exercised by Ministers of the Crown) 3. Constitutional conventions (accepted norms of political behaviour) 4. Common law (decisions by senior courts that are binding on lower courts) 5. EU Treaties 6. Statements made in books considered to have particular authority Note that not all of these sources form part of the law of the land, and so the British constitution encompasses a wider variety of rules, etc. than that of (say) the United States. 1. Statute law. Many important elements of the British constitution are to be found in Acts of Parliament. In contrast with many other countries, legislation affecting the constitution is not subject to any special procedure, and is passed using the same procedures as for ordinary legislation. The most important statute law still in force and affecting the constitution includes the following: 2. Royal prerogative. Certain powers pre-dating the establishment of the present parliamentary system are still formally retained by the Queen. In practice almost all of these powers are exercised only on the decision of Ministers of the Crown (the Cabinet). These powers, known as the "royal prerogative", include the following: 3. Constitutional conventions. Conventions are customs that operate as rules considered to bind the actions of the Queen or the Government. Conventions are not part of the law, but nevertheless are often considered to be just as fundamental to the structure and working of the constitution as the contents of any statute. Indeed, statute law affecting the constitution is often written in such a way that the existence of certain conventions is taken for granted, and some conventions are so fundamental that many people are unaware that they are in fact "unwritten" rules. Examples of the more important constitutional conventions include: 4. Common law. The common law is that part of the law which does not rest on statute. Instead, it is the accumulation of specific judicial decisions set by senior courts as precedents binding on lesser courts. Certain parts of the common law are also what is known as "trite law": examples of this include the fact that the United Kingdom is a monarchy, and the fact that brothers used to take precedence over sisters in the succession to the throne; this was repealed by the Succession to the Crown Act 2013. 5. EU Treaties. As a member state of the European Union, the United Kingdom is bound by EU law. However, the majority of votes cast favoured the referendum to exit the EU. The UK invoked Article 50 of the Treaty on European Union on 29 March 2017, which started what is often referred to as the "Brexit". The United Kingdom left the European Union on 30 January 2020. 6. Authoritative statements. Certain published works are usually considered to have particular authority. In the first half of the twentieth century this was the case with A V Dicey's "Law of the Constitution", being cited with approval in judicial decisions. A particularly important work is "Erskine May", which sets out the procedures and customs of the House of Commons. Other important sources include certain ministerial statements. However, none of these works have legal authority; at best, they are merely "persuasive".

UK Constitution and Government/Devolved Administrations. Devolution. "Devolution" refers to the transfer of administrative, executive, or legislative authority to new institutions operating only within a defined part of the United Kingdom. Devolved institutions have been created for Scotland, Northern Ireland, and Wales. Devolution differs from federalism in formally being a unilateral process that can be reversed at will; formal sovereignty is still retained at the centre. Thus, while the US Congress cannot reduce the powers of a state legislature, Parliament has the legal capacity to even go so far as to abolish the devolved legislatures. Devolution in Wales was originally restricted to the executive/administrative sphere, whereas in Scotland and Northern Ireland devolution extended to wide powers to pass laws. Scotland. The Scottish legislative authority is the Scottish Parliament. The Scottish Parliament is a unicameral body composed of 129 members (called Members of Scottish Parliament, or MSPs) elected for fixed four-year terms. Each of 73 members is elected by a constituency. The remaining are elected by eight regions, with each region electing seven members. Each voter has one constituency vote—cast for a single individual—and one regional vote—cast either for a party or for an independent candidate. Regional members are allocated in such a way as to permit a party's share of the regional vote to be proportional to its share of seats in the Scottish Parliament. The Scottish Government is the executive authority of Scotland; it is led by the First Minister. Other members of the Scottish Cabinet are generally given the title of Minister. The First Minister must retain the confidence of the Scottish Parliament to remain in power. Scotland has responsibility over several major areas, including taxation, criminal justice, health, education, transport, the environment, sport, culture and local government. The Parliament at Westminster, however, retains authority over a certain number of "reserved matters". Reserved matters include foreign affairs, defence, immigration, social security and welfare, employment, and general economic and fiscal policy. Wales. The National Assembly for Wales is the Welsh legislative authority. It is, like the Scottish Parliament, a unicameral body; it also uses a similar electoral system. Forty of its sixty members are chosen from single-member constituencies, while the remaining twenty are regional members. (There are five regions.) The Welsh Government is led by the First Minister and includes other Ministers, who must retain the confidence of the Assembly. The third Welsh Assembly can legislate using a system called "Assembly Measures". This system is a lower form of Primary Legislation similar to Acts of Parliament. They can be used to repeal laws, create provision and amend laws. The difference with "Assembly Measures" and "Acts of the Assembly" is that Measures do not have a bulk of powers with them, each Measure will come with a LCO, or Legislative Competency Order, which transfers powers from the UK Parliament to the Welsh Assembly Government. Devolution in Wales has changed a lot since 1999. In order for the National Assembly to have full legislative powers, they will need to trigger a referendum through both the Assembly and both houses of the United Kingdom parliament. Once done, Wales will for the first time ever, will be able to legislate and make their own Acts. (To be known as Acts of the Assembly, or Acts of the National Assembly for Wales). In early 2011, a referendum held in Wales approved the transfer of full legislative competence to the National Assembly in all devolved matters. Northern Ireland. Northern Ireland was the first part of the United Kingdom to gain devolution, in 1921. However, it has had a troubled history since then, caused by conflict between the main "Unionist" and "Nationalist" communities. Because of this historical background, the present system of devolution requires power to be shared between political parties representing the different communities, and there are complex procedural checks in place to ensure cross-community support for legislation and executive action. The Northern Ireland Assembly comprises 108 members elected to represent 18 six-member constituencies. The Executive (government) is made up of members from the largest parties in the Assembly, with ministerial portfolios allocated in proportion to party strengths. The Executive is headed jointly by a First Minister and Deputy First Minister, who are jointly elected by the Assembly. The Assembly's legislative powers are broad, and are similar to those of the Scottish Parliament (with the notable exception of taxation). The transfer from the United Kingdom's central government of responsibility for the criminal justice system has been highly contentious, and has only recently been carried out.

Electronics/Microwave. Microwave. Microwaves are defined to be waves having frequency greater than 1 GHz. While some books put the upper limit of 30 GHz to microwaves, some literature consider the upper limit to be up to 300 GHz. At such high frequencies, the behaviour of the devices used commonly is no longer predictable by general rules and new design issues need to be considered. Even short leads of devices act as inductances and may possess capacitances. A magnetron emits microwaves at 2450 MHz λ = 12.24 cm A microwave is a Faraday cage, so waves can't escape openings smaller than 6.12 cm. Microwaves work by bombarding atoms from all sides with an electric field. This causes the dipole to keep realigning itself and the food heats up. Most effective on water, metal deflects microwaves. Containers survive because they do not absorb many microwaves relative to the food.

Electronics/Photocopier. The Photocopier. One example of the practical use of static electricity is a photocopier. A photocopier is a fairly complicated piece of equipment, but the basic principle of how it works is fairly simple. The best way to understand what is going on is to consider it as a stage by stage process. Stage one. Positive charge is sprayed onto a plate from a high voltage power supply. The plate is connected to the earth but the charge does not have quite enough energy to flow away. (The plate is not a good conductor of electricity) Stage two. Paper is placed over the plate and a light shone onto the paper. Where the paper is white the light is reflected onto the plate, where the paper is dark a shadow falls onto the plate. The light falling on the plate gives it just the extra energy needed to allow the charge to escape to earth. The plate becomes neutral where the paper was white but keeps its charge where the paper is black. The plate is now a copy of the paper with charges taking the place of ink. Stage three. Toner particles are spayed through a negatively charged nozzle onto the plate. As the toner passes through the nozzle it picks up the charge so that each particle of toner becomes negatively charged. The now charged toner is attracted to the areas of positive charge because unlike charges attract. More light then allows the positive charge to escape (however, the negative charge on the toner remains). Stage four. A piece of paper is given a very strong positive charge, and then placed in contact with the plate. The paper attracts the toner. The paper is then removed from the plate and passed through a heating unit. The heat melts the toner and bonds it to the paper. In a real photocopier, there is no plate, just a large drum. As the drum rotates its surface goes through stages one to four. At the end of the sequence a scraper removes any toner left on the drum and the whole process is repeated with a new image. A good photocopier can process a page in less than a second. Questions on Photocopiers. 'Q1)'If the black areas of the image leave a positive charge on the plate what charge do the white areas of the image leave? ("Be very careful, the answer may not be what you think!") Q2) How does the light shining onto the charged plate allow it to lose its charge ? Q3) Why is the toner given a negative charge ? Q4) Why does the paper attract the toner ? Other uses of static electricity. Spray painting car parts. When paint is sprayed from a paint gun, the painter normally needs to use a fair amount of skill to ensure the paint goes on evenly. By connecting the spray nozzle to a negative electrode, it is possible to charge each droplet of paint. If the car part is then given the opposite charge, the paint droplets will be attracted to the car body part. This has several advantages. Questions. Q5) If the paint droplets are given a positive charge, what charge should the car part be given? Q6) Why does less paint fall on the floor?

Electronics/Devices. < Electronics/Expanded Edition Specific devices made through applying knowledge found in this book.

Intelligence Intensification. Authors. __NOEDITSECTION__

Intelligence Intensification/Introduction. The need of yet another book on this subject. All books on this subject I have come across have either one or both of the following disadvantages. Firstly, many books are filled with advertisements for useless products, workshops and so on. They are designed for personal development junkies who only want to feel the joy of reading books, listen to tapes and buy devices which "could have" helped them if only they were willing to work hard. The second disadvantage is that most books are thin on information—what could be explained in a couple of pages of text is explained in hundreds of pages. This book will be different. The information in the book is dense and only intended to inspire the reader's "own" thinking and experimentation. I suppose the author's desire for yet 'another book on this subject' is more a result of the need to define his topic than a dearth of information. A serious attempt at this endeavor would be better served by one with a little more education on the subject. An introduction of this topic more than most others, would benefit from defining your goal and organizing an approach. As with any discussion of intelligence, definitions, frames of reference, processes, development (natural or guided), the relationships between original, developed, derived, summarized, or compiled knowledges, cultural variation, and scales of valuing would be some primary considerations for interactive descriptions. Many references exist for each of these subtopics, and one anticipated problem will be debates on defining intelligence first, then (and possibly unnecessarily) discussion on measuring. Post-graduate (as fair a starting point as another) studies of assessment deal with the evolution of the qualitative/quantitative measuring of intelligence (and primarily within educational or legal/social contexts). All of this would be logical to establish a platform from which one explores methods of intensification or development of intelligence.-- () 17:24, 21 July 2008 (UTC) I envision two directions for an introduction, one evolutionary and the other categorical, with a strong allowances for aristotelean vs. oriental processes of perceptualization/development of understanding. -- () 17:24, 21 July 2008 (UTC) Binet, Gardner, Cut the Crap. In keeping with the aim described in the preceding paragraph, this book will not contain affirmations that these techniques will work "even for you" and so on. No unscientific testimonials will be included. In a world with a lot of information it is of utmost importance that the information is compact, concise and to the point. Spread the Word. The title "Intelligence Intensification" establishes that the goal of this text is to intensify intelligent thought in the reader. Perhaps more importantly, the book aims at the greater goal of increasing intelligence all over the world—for joy and for the benefit of all humanity. This book is inspired by Timothy Leary's 20th ­Century concept: "Every citizen a Scientist". What's your vision? The following is a quotation from the book "The Art of Doing Science and Engineering: Learning to Learn" by Richard Hamming. He models a life with or without a vision with random walks. Before or while reading this book and practicing its contents, try to grow a vision of your life.

Intelligence Intensification/Memory Techniques. Visualization. Have you ever seen or heard of a person with a super memory? Did he remember all the cards in a deck in perfect order? Maybe he claimed that he had a photographic memory. That, however, is probably not the case. Rather, he or she used simple techniques which you can master too, with little practice (say, 100 hours). After half an hour of practice, you will see fantastic results. The memorizing time is directly related to the number of things to memorize. If you need five minutes in order to remember ten foreign words, you'll need ten minutes in order to remember twenty. It should be stressed that you will not get a photographic memory. You will be able to remember faces, foreign words, telephone numbers and much more. But you will not be able to look at a picture and then immediately remember every detail in your inner vision. Search the web for books on memory improvement and you will find lots of books claiming that they will give you a photographic memory. This is generally a lie. The systems you will learn on this page, however, are extremely practical and simple to use. And best of all, they work! The basic theory is that your brain associates everything you learn with what it already knows. Using the memory techniques to follow, you consciously make those associations (in addition to the ones made unconsciously). The other cornerstone used in the techniques is that the brain remembers strange things better than mundane things. This is the absolute essence of memory techniques. Lets see what comes out of it! Remembering a list of items. The two basic methods used for remembering a list are linking and pegwords Linking. After visualizing the first item on your list, visualize it associated with the next item, then visualize that item associated with the third, and so on. For example, to remember a list like "apple, fish, lady, star, stop sign, pencil ..." imagine an apple. Now you shall link this apple with a fish by visualizing (for instance) an apple tree with fishes instead of apples-you could even imagine a fish falling on Newton's head or Eve handing Adam a fish. Remember that it should be weird. Next link the fish with lady by visualizing a mermaid. Next item: Visualize a night sky with shining ladies in the sky instead of stars. Link to stop sign by visualizing a falling star landing on the ground-only instead of a star, when you get up close it's a stop sign, link "stop sign" with "pencil" by imagining a stop sign which is held up not by a metal post, but by a giant pencil. Pegwords. Pegwords are used if remembering the position of each item is helpful; for example, if you would like to be able to recall the fifth U.S. president. In this system, a list of pegwords is pre-memorized. The pegwords are designed to be easy to visualize, and to associate with a number. For example memorize this list of pegwords. Try to form a mental image of each one. And then to remember any list of 10 items you associate the items to be remembered with their pegs. To memorize the grocery list: apples, butter, razor blades, soap, bread, milk, cat food, bacon, batteries, and orange juice, Sound obscure? Close your eyes and try to remember the list. Writing The List Down. Sometimes, writing down something you wish to remember may help to increase your retention. If you have the time and resources available, you may write down the entire list, maybe a few times, and be able to remember most, if not all of the items. You may also use this to remember paragraphs of text as needed. Method of Loci. The Method of loci was used by the ancient Greeks for oration; that is, for memorizing speeches. It involves mentally transversing a familiar place, for them, a temple, stopping at certain points, and associating those points with sections of a speech. /Memorizing Numbers and Digits/. A system known as the mnemonic major system used to convert numbers into words. See module for details. Techniques Which Combine Pegwords With the Number System. /A Longer Peglist/. While rhyming peglists are simple to learn they are severely limited in the number of pegs that can be created. 7 rhymes with 11, 21 rhymes with 31 rhymes with 41, etc. Using the major system, rather than rhymes, you can create as many pegs as you want with no ambiguity. This module currently has 116 pegs- enough to accommodate the periodic table. /Memorizing Playing Cards/. This module combines the major system with the pegword systems to do some pretty amazing things with cards. It does require a lot of practice, and requires you to be adept at both systems, but you will be surprised with what you can accomplish. /The Periodic Table/. An example of just how powerful this system can be. Pegs for the atomic numbers linked to a pun on the element names or their symbols linked to a phrase describing the atomic mass. This module has yet to be written. You can help. Remembering foreign words. In order to remember a foreign word, or a native word you don't know, all you need to do is to make it concrete in your native language, and then link it with its meaning (which maybe has to be made concrete too). To remember the French word for horse, cheval, make it concrete in english by thinking of the word chew (which is a concrete verb). Then link a horse with chew. Maybe you vividly visualize that you are watching a big, big horse chew some trees from the comfort of your hammock. Then he picks up the trees your hammock is swinging between. Your blood splashes in your face as your bones crunch and you are racked with pain. You curse the horse for chewing you up. So when you think of horse, you immediately think chew, which reminds you of cheval. In order to remember the Swedish word for luck, tur, you also need to make the native word luck concrete. One easy way is to think of winning a lottery. But emphasize the middle syllable of "lottery": lot-TUR-ry. Think about how much TUR you need to win the lot-TUR-ry. Remembering faces, art works etc. This is easy if you know any of the other techniques. To connect a face with a name, pick out something in the face and link whatever you picked out with the name. If you meet a girl whose face is shaped like a heart, link a heart with her name. So if her name is Angela, you can imagine an angel descending from heaven, who violently rips out her heart. Imagine her screaming, the blood spraying everywhere, and her squishy heart still beating while falling down to the muddy earth. If you meet somebody with a funny nose, it's easy. If you don't find anything special (but isn't everyone special in some way?) just pick something. What do the lips look like? The ears? The total impression? The technique to remember the author of a painting is analogous. Only by "trying" to find something special you are forced to really look, which helps remembering. How to remember whatever you want. Combine the above and below techniques, use them as a basis to construct your own personal techniques and remember the essence: Use weird, vivid visualizations; use all your senses (taste, smell, etc); associate what you want to remember with as much as possible. See Wikipedia's list of mnemonics. Drugs and Vitamins. Some diet supplements, vitamins, herbs and amino acids may also improve your memory and alertness, though very few if any have been scientifically studied for effects on memory. These include the B vitamins, caffeine, amino acids like phenylalanine, (though an overdose can be deadly) melatonin, ginseng, ginkgo biloba, and many others. There are also a number of prescription-strength drugs which might be legal, but may have side-effects although current studies show some of them to be safe to use. Some of these alleged "smart drugs" include Vasopressin, Hydergine, and Piracetam, and DHEA. You should consult your physician before taking any drug or supplement, especially if you are taking any other medications. Current research suggests the best ways to keep your memory healthy as you age are to use it frequently (read non-fiction, engage in political conversation, etc.), eat a well-balanced diet, take antioxidant supplements (though overdoses of selenium and vitamins A and E could be fatal), and possibly increase the proportion of omega-3 fatty acids in your diet.

Intelligence Intensification/The Concept of Change. Speculative History of the West. For thousands of years there was little change toward better conditions for human beings. The changes in the ability to treat diseases was minor. Technology improvement was slow. Then, suddenly, rapid change occurred. Initially, people didn't expect change. As a self-fulfilling prophecy, nothing happened. On many levels change was not tolerated. Then, new inventions became common, change was commonplace and people began to expect change. And what was expected, happened. And it happened even more dramatically than anyone expected. Today, new knowledge accumulates and the ability to apply it is acquired faster than ever. The Age of Networking is born. Also, during the centuries, old systems have been replaced by newer more efficient systems. A new government replace the old with each new revolution. After a while, the new system is old and it's time for a new replacement. On a Smaller Scale. On a smaller scale, this can happen in an individual too. Use the power of self-fulfilling prophecies. Love change and don't be afraid of changing. Every time we learn something, we change. If you are afraid of changing it will make learning harder. Moreover, each remodeling of the solution to a problem tends to make the solution better. Everyone who has done some computer programming knows that a program solution tends to be better if it is remade from scratch instead of applying ad hoc patches. Use this phenomena to your advantage. Throw all your habits out of the window every now and then. Start from scratch, use your new and increased knowledge to create better habits. Our nervous system receives a large amount of impressions each second. From this chaos it selects a pattern which we perceive as reality. Thus our perception of reality is colored by how we recognize patterns. We see the patterns we know. Example: Think of red. Look around right where you are sitting now. Your brain extracts everything that is red. But what smells were around? How many green items did you see? The brains works in the same fashion when it comes to ideas, views and opinions. Thus, if you try to build on your old knowledge new information will be filtered according to your old information. A better idea might be to try to start from scratch, as described above.

Electronics/Expanded Edition Fundamentals. ../Expanded Edition/ Coulomb's law. Dielectric constant, definition of charge Electricity The force resulting from two nearby charges is equal to k times charge one times charge two divided by the square of the distance between the charges. The electric field created by a charge is equal to the force generated divided by the charge. Electric field is equal to a constant, “k”, times the charge divided by the square of the distance between the charge and the point in question. Electric potential energy is equal to a constant, “k” multiplied by the two charges and divided by the distance between the charges. Variables F: Force (N) k: a constant, 99 (N•m2/C2) q1: charge one (C) q2: charge two (C) r: distance between the two charges, (m) A charge in an electrical field feels a force. The charge is not a vector, but force is a vector, and so is the electric field. If a charge is positive, then force and the electric field point in the same direction. If the charge is negative, then the electric field and force vectors point in opposite directions. A point charge in space causes an electric field. The field is stronger closer to the point and weaker farther away. Ampere's law. Definition of current, magnetic permittivity

Electronics/Expanded Edition DC circuits. < ../Expanded Edition/ Resistance. Ohm's law, resistivity and conductivity Capacitance. Some formulae for capacitance of different shapes, dependence on dielectric constant, units Batteries. Behaviour in circuits Solving circuits. Kirchoff's laws, Thévenin's theorem Hall effect, Piezoelectricity, etc. Voltages induced by magnetic fields, stress, temperature difference, etc.

Invertebrate Zoology/Insects. Subclass Apterygota Subclass Archaeognatha (jumping bristletails) Subclass Dicondylia Subclass Introduction to the Insects. In terms of number of species described, the Subphylum Hexapoda is the largest in the Phylum Arthropoda. Hexapods (meaning "six legs or feet") include the insects (Class Insecta), which dominate terrestrial environments on earth. The subphylum also includes several groups of wingless arthropods of uncertain placement; namely the orders: Diplura, Collembola, and Protura.

Electronics/Expanded Edition Impedance. Inductance. Mutual and self, definition, dependence on geometry Impedance. Combining R,L, and C into one quantity. Solving the circuit. How to set up and solve the ODE's, free and driven circuits.

Electronics/Expanded Edition Resonance. < ../Expanded Edition/ Resonance. Simple resonant circuit, description. Filters. Ideal Filters. This section introduces first order butterworth low pass and high pass filters. An understanding of Laplace Transforms or at least Laplace Transforms of capacitors, inductors and resistors. Low Pass. Transforming the Resistor and Capacitor to the Laplace domain we get: Expressing formula_3 using terms of formula_4. The transfer function is So For the Frequency Domain we put formula_9 The magnitude is and the angle is As formula_13 increases formula_14 decreases so this circuit must represent low pass filter. Using the -3 dB definition of band width. formula_15 Therefore Which gives the general form of a low pass butterworth filter as: , where k is the order of the filter and formula_20 is the cut-off frequency. High Pass. If all the component of the circuit are transformed into the Laplace Domain. The resistor becomes formula_21 and the inductor becomes formula_22. Using voltage divider rule formula_23 below is reached. If formula_23 is transformed into the frequency domain by putting formula_26. Which has a magnitude of and an angle of

Electronics/Expanded Edition Aerials. This is about what aerials and antennae look like to a circuit, effective impedance, coupling of current with radio waves. Intro to design of radio receivers/transmitters, including much maths I. Introduction to Electromagnetics / Transmission Line Theory Previous sections have presented an overview of DC and AC circuit principles. The circuits we have looked at until now all have one thing in common: they are "electrically small" compared to the "wavelength" of the voltages and currents induced on them. In other words, the "physical" size of the circuit is much smaller than the length of the voltage or current wave. Many other circuits, however, have physical dimensions the same as or larger than the wavelengths of their voltages and currents. In order to understand their operation, it is necessary to introduce the principles of "electromagnetic" theory. All alternating voltages and currents move through space the form of waves. This is a fundamental principle of physics discovered in the 19th century and elaborated by the Scottish physicist James Clerk Maxwell, among others. Like a wave in water, an electromagnetic wave contains "crests" and "troughs". "Wavelength" is the distance from one crest to the next. In the electrical case, a crest is a point in space where the alternating voltage or current value is at maximum, and a trough is a point where the value is at minimum. In the case of a DC current or voltage, the value is the same everywhere, so the wavelength is infinite. Alternating voltages and currents have finite wavelengths. The pattern the values take in space is that of a "sine wave". If voltages and currents are measured at locations along the path of an electromagnetic wave, they will fit the curve of a sine wave. "Frequency" in electromagnetics is the number of crests passing a particular point in space per unit of time, as the wave moves. It is the "inverse" (one over) wavelength. What has been discussed so far is a "subset" of "electromagnetic" theory.

Electronics/Expanded Edition Afterword. Electrolysis. Description, Faraday's rule, battery in reverse, maybe fuel cells Ultra high frequency AC. Show how things change at ultra high frequency, wires becoming waveguides. This points the way to unification of electricity and radio. Maxwell. Math free description, ubiquity of EM Quantum stuff. Math free Computers. Path from small circuits to PCs Final word. Math free piece on how electronics fits into the wider scheme

Intelligence Intensification/Creativity Techniques. The techniques described here are not very useful for a musician or a painter (or they are, but not directly). The goal is to get new ideas and break old thinking habits. The essence of the theory is as follows: Usually your mind is in a rut. In order to move your thinking wagons out of the rut, you provoke it. You will find it to be very simple and it will work immediately. With practice, it will work even better! (Here it would be nice with a picture with thinking wagons stuck in a rut. Can you draw one?) Sometimes the brains seems to be slow, lacking energy and wit. Maybe you have to finish an article before you go to bed. Ideas don't flow easily. To get the brains into full swing so thinking goes with a bang, you might try the following exercise: Point at an object in your surroundings and call it by another name. Point at a lamp, and say snake! Point at a floppy disk, say spoon! A book -- dance! Table -- Swimming! Etc. Try to increase the speed until you cries out weird words immediately when you point at something. I will not say how it feels afterwards. Do it and you will know! You might observe that it is easy to slip into patterns, like "Cow! Milk! Breakfast! ...". Practice a little bit more to get rid of these patterns. Now I will give you two good techniques to get new ideas. I think it's all you need. There are thick books about this, and you may read them if you think it is interesting. But if you are intelligent, understanding the essence is all you need (everywhere, always, in every theory. Am I overenthusiastic?) The first one is to simply have something to think of. Instead of looking at the whole problem at once (if you have a problem), look at a tiny detail. If you walk down the street, force yourself to only look at -- for instance -- a lamppost or maybe even only at the hold or the very lamp. Then think, how could this be done better? Compare this to looking at the whole street at once, thinking "Can anything be done better here?". Make this a habit if you wish to be the guy or girl who always have new and funny ideas. It is like a burning glass which produces a warmer spot if the spot is smaller. So have a small focus point! The second thing, which you can combine with the above idea to get good results is to have a thought in mind, a thought you want to develope. Then provoke your brain with a random word, like "pancake batter!" or "machete!". Then let this word bring new ideas. Often all you need is to be provoked. The very provocation, the weird word, is only needed to bring the thinking out of the rut. Here is an example: In my hands I've got a CD case. I want to make it better. As a focus point, I choose the hold. I need a provocation word. FROG! Frog? They are slimy... Gluing the disk is not a good idea... But one could hold the disk with suction cups! There is actually free space on the plastic of a CD where no data is written, so leaving suction marks doesn't matter. Maybe this is a very stupid example, but this is how it works even when you're thinking of not so stupid things. Remember the essence! You must provoke your brains. People who haven't provoked their brains in years get dull and robotic. People who provoke their brains often get creative and non-robotic. Additional Techniques for Provoking Creativity and Intelligence Random associations. The point of this exercise is to be able to state how a ___________ is like a ________________. Use whatever method you choose to fill in the blanks and state how the two are alike. For example, suppose you want to make a game of this with a child. Place assorted objects into a bag. Gather the objects from various places, so that they are not all related by location, toy figures count as well. Then have the child, without looking and, doing so quickly as not to select by feel, pull out two objects. Then have the child describe how the two objects, say a bottle opener and a (toy) giraffe are alike. Additional ways of filling in the blank might be to randomly select words from a dictionary of telephone yellow pages, select images from advertisements or commercials (it takes some of the mindlessness out of television watching), with multiple people have each select a word individually secretly (nouns work best) then fill in the blank with the two word and see who can come up with the most connections. No-name naming. Describe a scene, or the room that you are in without using nouns, except for the names of geometric shapes. Use only adjectives and verbs. Try describing a specific picture from a group of pictures that are not too dissimilar with enough detail that another person can pick out the one you are describing. Smallest differences. From a group of 4-6 similar objects, such as lemons, oranges, eggs, unsharpened pencils, choose one specimen. Spend 5-10 minutes getting to know your specimen. Then return it to the group and mix them up. Test yourself to see if you can differentiate yours from the rest. (This is a useful skill when doing magic tricks with cards. If you can discern small differences, every deck of cards becomes a marked deck with use.) Observation. This is the most basic skill of all, and the one others are based upon, the ability to observe. Where ever you are, take time to observe and verbalize (even if only internally) what it is you observe. Start with one sensory channel, such as visual, then proceed to auditory, then kinesthetic, and olfactory. Looking at things upside-down. Or sideways. Turn your head so that you look at things upside down. Think about what would happen if suddenly gravity changed directions. Or imagine the building you are in suddenly turned over on its side and imagine how you would climb out. Sudden difference. Look at something and imagine what it would be like if a particular detail about it changed to its opposite. For instance, what if CDs were suddenly square? Different use. Think of alternate ways to use things than for the original purpose for which they were intended. For example, children often consider their parent's bed a trampoline instead of a bed. Kitchen sink could be used as a mixing bowl. For that matter, mixing bowls could be made into sinks! CDs could be used as frizbees.. or light catchers... or christmas tree ornaments. I've heard of people using doors as desk surfaces. (This kind of exercise is particularly helpful for conserving money in times of unemployment.) Supplement: Creativity is simply making new intellectual connections. Everyone has some potential to be more creative than they are now; it just takes a willingness to allow their minds to wander in an unforeseen direction. When an unusual thought enters many people's heads, they dismiss it as silly or pointless. They may be right, but that's beside the point. By judging these spurs of creativity before the mind is given a chance to follow them to their conclusion, they are closing themselves from creative thoughts that may be relevant, useful, or entertaining at the very least. Next time you get an odd idea, run with it. If you see a man walking his dog in the park and you begin to wonder how things would be if the dog were instead walking the man, follow that thought and see where it takes you. It may sound silly, but sharpening your creative process for these simple observations will sharpen your ability to come up with original solutions to things more relevant to your life too. Be creative! You have nothing to lose and only a sense of humor, a newfound ability to express yourself, and an ability to solve problems better to gain.

Intelligence Intensification/The Theoretician Versus the Practician. The Theoretician. In periods I'm a striking theoretician. I think about a lot of interesting things and new ideas pops up in my mind all the time. But nothing gets done. That's the problem of the theoretician. Knowing that one is a typical theoretician, one has better chances to counteract. The typical situation is that one buys a book or reads a text about something interesting. Instead of putting everything into practice one contents oneself with reading the exercises (if there are any) and thinking about the content. One may actually feel that one really learnt something and that the reading gave lots of insights. Imagine how many insights you would get if you tried everything you read about in physical reality (whatever that is). Another problem associated with books is that authors often wants to write hundreds of pages about a theory which very well could be explained in ten pages. In general, all an intelligent reader needs is the essence of a theory. Everything else is evident when the essence is understood. In fact, many books explains so many unnecessary details that the essence is never understood. That is the problem of the theoretician. He spends time reading about details which soon becomes obvious for anyone who put the ideas to use. Whenever you, the reader, feels a section in this book doesn't break down to the essence of the subject, please use the possiblity to discuss a page in order to make the active contributors of this book aware of the problem. The Virtue of Discipline. The only way to solve the problem of putting things into practice that I can think of is discipline. Maintaining discipline is crucial when one must overcome the psychological barrier that so often appears when one wants to learn something or do something creative. Discipline can be practiced in various ways, the reader can probably invent his own methods. You are your best teacher. That is why you should only read this book to inspire your own thinking and practice. ???formula_1

Intelligence Intensification/Repetition. Why Boring Repetition? "If you want to" minimize studying time, repetition is important. If you repeat, you tell your brain this is important, remember this. Optimizing Repetition. "Repetition should be made" at certain intervals of time. Of course it's best if you experiment and find your own optimal intervals but a general guideline is as follows: After one hour of study wait ten minutes. During these ten minutes your brain will organize the new information. Then you should briefly repeat the information during about five minutes. The next repetition should be after about an hour, also about five minutes long. Then repeat again after one day, as much as you feel needed. Try again after a week and then after a month. Hopefully the information is in your long-term memory by then. Note that these are only general guidelines which may very well vary due to the nature of the material you study, your studying technique, etc. Experiment! Aiding Software. There already exist several software programs that optimise the repetition cycle. SuperMemo(proprietary) uses a repetition algorithm, which is based on neural research. Free programs with a similar algorithm include such programs as Anki or Mnemosyne. The best combination between using a structured approach and optimal repetition intervals is realised with RecallPlus(proprietary). RecallPlus uses concept maps for structuring the learing material and then uses repetition cycles, associated with the links between topics.

Intelligence Intensification/Speed Reading. Classical Speed Reading. If you have never thought about how you can increase your reading speed, you will find these tips very interesting. You can probably read at much higher rates without losing understanding. Your brain is faster than your eyes. If you're still skeptical about increasing your reading speed, consider the following. This technique does not have to be used to increase your reading speed. You could choose to continue reading at your old speed, but have more time to understand the text, with less time used interpreting symbols. It is also strange that many people say that they don't need to read faster, but nobody wants to read slower. Speed reading is not suitable for many types of written material. Text with complex mathematical concepts or equations, for example, needs to be read many times, slowly. Written material describing a complex historical scenario likewise may be too complicated to be understood quickly. However, when the reading speed is slow and the information is not interesting, the mind tends to wander. Speed reading poetry will not improve one's appreciation of it, partly because in reading it too quickly, one loses much of the nuance of poetical metre. This is also true if you try to recite poetry at 700wpm. Speed reading tends to be a useful technique for light material such as easy-reading novels or magazine columns. The best idea is to speed read slowly. That is, use a perfect reading technique, but do it at a rate appropriate for the appreciation of what you are reading. Using a bad technique doesn't increase the appreciation of poetry, or the understanding of mathematical formulae. Speed reading the small print of speed reading course contracts (or any contract, for that matter) is also not a good idea. In order to increase your reading speed, do this: Aiding Software. The key to success is concentration and therefore discipline. If you really want to practice the last point, reading as large chunks of words as possible, you can create a computer program which writes two words on the screen with a space between them. Both words should be read at once, not one after the other. As you practice, enlarge the space between the two words. Another idea is to write a program which shows a text and as you press the space bar it begins to make the text invisible one word at a time (or maybe x words at a time). In that way, you are forced to continue reading, you cannot look back. If you write any of these programs, which shouldn't take more than a couple of hours, please make them available for everyone here. ("Shouldn't the code be posted to , under ?") Yet another option which may be more viable for some is using available speed reading software package. Among the best are The Reader's Edge, RocketReader and AceReader, though of course many others exist. Most of these programs will have grouping, flash, memory, speed reading, and shadow exercises.

Intelligence Intensification/General Tips For Practicing. Information Ecology Practices. The transition to a knowledge-based digital environment has been vital to the transformation of the nature of information ecosystems and a global dynamic of the intensification of intelligence. The emerging holistic life science of Information ecology incorporates principles, disciplines and protocols that facilitate the emergence of a common collective intelligence: Doing One Thing At a Time. In karate, practicing the basic strike demands that Having all this in mind when practicing is very difficult. A better idea is to concentrate on one thing at a time. Iterate over the list of things to practice. Let the other things happen by themselves or go wrong, as the case may be. This method of practicing is of course applicable in all kinds of training. Choosing to Focus on the Practice. Another useful technique to use when practicing, also inspired by martial arts, is to enter a meditation position before and after the training. Before the training, enter the position and After the training, enter the position and

Electronics/Vacuum Tubes. Background. The first vacuum tube diodes were created by Thomas Edison in 1904. These initial tubes could only be used for rectification. The triode, which allowed for voltage and power amplification, was invented 3 years later in 1907 by Lee De Forest. Vacuum tubes are also referred to as "thermionic tubes", "thermionic valves", "electron tubes", "valves" and just plain "tubes". While, for most electronics, vacuum tubes have been replaced by transistors, there are still some uses for which vacuum tubes are desired. Vacuum tubes are frequently used in high end Hi Fi amplifiers and are generally desirable over transistors for their "warmer" tone. They are also generally preferred in guitar amplifiers both for their smoother clipping in overdrive and the warmer tone. Finally, vacuum tubes are used in high frequency communications (at frequencies that would destroy solid state components) as well as satellite and military communications due to their durability (they stand up to solar radiation better than solid state and are immune to electromagnetic pulse). Passive Versus Active Components. Passive components have no gain and are not valves Vacuum Tube Basics. A Vacuum Tube is a container (usually of glass) from which the air is removed. Inside the tube are two or more "Elements". - Cathode: (electron emitter) has an electrically heated filament (which you can usually see glow red ) which spits out electrons that travel through the vacuum to the Anode (electron acceptor). - Anode (a.k.a. Plate): Is a conductive (usually metal) plate that is connected to a positive voltage. The negative electrons flow from the Cathode to the Anode. A vacuum tube with just Cathode and Anode elements is a DIODE. Current will flow only when the Anode has a positive voltage relative to the Cathode. - Grid(s): metal gratings or grids are placed between the Cathode and Anode to produce devices that can amplify signals. NOTE: A tube with 3 elements (one grid) is a TRIODE, with 2 grids a TETRODE, with 3 grids a PENTODE. A Grid between the cathode and anode controls the flow of electrons. By applying a negative voltage to the grid it is possible to control the flow of electrons. This is the basis of the Vacuum Tube amplifier. Vacuum tubes contain heater filaments. These are similar to the filaments one would find in a standard light bulb. The filaments usually run at low voltages (6V and 12V are common, though tubes can be found with a variety of filament voltages). The filament is surrounded by a cathode. The filaments are constructed of heat the cathode to about 800 degree Celsius, at which point the cathode begins to emit electrons. The electrons normally float on the surface of the cathode and this is called a "space charge". The anode is generally kept more positive than the cathode and so it draws the electrons off of the cathode. The voltage difference in the direction from the cathode to the anode is known as the forward bias and is the normal operating mode. If the voltage applied to the Anode becomes negative relative to the Cathode, no electrons will flow. In electronics vacuum electron tube or valve is a device that controls current. Through a vacuum in sealed container. Vacuum tubes rely on thermionic emission of electrons from hot filament. Tube Characteristics. Below is an explanation of various tube parameters: Diodes. Vacuum tube diodes contain only two electrodes (besides the heater): A cathode and an anode. The filament heats up the cathode, producing a space charge. A relatively positive voltage on the anode then draws electrons from the cathode to the anode, producing the one way current of a diode. Current will not flow from anode to cathode. As with silicon diodes, vacuum tube diodes can be used for various functions. They can be used in voltage multipliers, envelope detectors, and rectifiers, for example. Common rectifier tubes are: 5AR4/GZ-34, 5V4-GA, GZ37, 5U4-G/GA/GB, 5Y3-G/GA, 5R4GYB, and 5R4-G/GY/GYA. Different rectifier tubes have different maximum voltages, current ratings and forward voltage drops. Some have 5V filaments and others 6.3V. And some draw more filament current than others. Some rectifiers are half-wave (single diode) and some are full-wave, containing a single cathode, but two anodes, one for each half of the wave, as shown in the Full-wave rectifier image to the right. Klystron. A klystron is a vacuum tube used for production of microwave energy. This device is related to but not the same as a magnetron. The klystron was invented after the magnetron. Klystrons work using a principle known as velocity modulation. The klystron is a long narrow vacuum tube. There is an electron gun (heater, cathode, beam former) at one end and an anode at the other. In between is a series of donut shaped resonant cavity structures positioned so that the electron beam passes through the hole. The first and last of the resonant cavities are electrically wired together. At the cathode the electron beam is relatively smooth. There are natural slight increases and decreases in the electron density of the beam. As the beam passes through the holes of the resonant cavities, any changes in the electron beam cause some changes in the resting electro magnetic (EM) field of the cavities. The EM fields of the cavities begin to oscillate. The oscillating EM field of the cavities then has an effect on the electrons passing through, either slowing down or speeding up their passage. As electrons are affected by the EM field of the first cavity they change their speed. This change in speed is called velocity modulation. By the time the electrons arrive at the last cavity there are definite groups in the beam. The groups interact strongly with the last cavity causing it to oscillate in a more pronounced way. Some of the last cavity's energy is tapped off and fed back to the first cavity to increase its oscillations. The stronger first cavity oscillations produce even stronger grouping of the electrons in the beam causing stronger oscillations in that last cavity and so on. This is positive feedback. The output microwave energy is tapped of for use in high power microwave devices such as long range primary RADAR systems. The klystron is a coherent microwave source in that it is possible to produce an output with a constant phase. This is a useful attribute when combined with signal processing to measure RADAR target attributes like Doppler shift. Related microwave vacuum tubes are the Travelling Wave Tube (TWT) and the Travelling Wave Amplifier (TWA). A Hybrid device, which combines some aspects of these devices and the klystron, is a device called a Twicetron. Magnetron. Magnetrons are used to produce microwaves. This is the original device used for production of microwaves and was invented during the Second World War for use in RADAR equipment. Magnetrons work using a principle known as velocity modulation. A circular chamber, containing the cathode, is surrounded by and connected to a number of resonant cavities. The walls of the chamber are the anode. The cavity dimensions determine the frequency of the output signal. A strong magnetic field is passed through the chamber, produced by a powerful magnet. The cathode is similar to most thermionic valves, except heavy, rigid construction is necessary for power levels used in most Magnetrons. Early, experimental designs used directly heated cathodes. Modern, high powered designs use a rigid, tubular cathode enclosing a heater element. Naturally excited electrons on the surface of the cathode are drawn off, into the chamber, toward the outer walls or anode. As the electrons move out, they pass through a magnetic field that produces a force perpendicular to the direction of motion and direction of the magnetic field. The faster the electrons move, the more sideways force is produced. The result is that the electrons rotate around the central cathode as they move toward the outside of the chamber. As electrons move past the entrances to the resonant cavities a disturbance is made to the electro magnetic (EM) field that is at rest in the cavities. The cavity begins to oscillate. When another electron moves past the cavities, it also interacts with the internal EM field. The motion of the electron can be slowed or sped up by the cavity field. As more electrons interact with the cavity EM fields, the internal cavity oscillations increase and the effect on the passing electrons is more pronounced. Eventually, bands of electrons rotating together develop within the central chamber. Any electrons that fall behind a band are given a kick by the resonant cavity fields. Any electrons going too fast have their excess energy absorbed by the cavities. This is the velocity modulation effect. The frequency of the resonance and electron interaction is in the order of GHz. (10^9 cycles per second) In order to have a signal output from the magnetron, one of the cavities is taped with a slot or a probe to direct energy out into a waveguide for distribution. Magnetrons for RADARs are pulsed with short duration and high current. Magnetrons for microwave ovens are driven with a continuous lower current. The magnetrons for WWII bombers, operated by the RAF, were sometimes taped into a shielded box so that the aircrew could heat their in-flight meals, hence the first microwave ovens. Cathode Ray Tube. A cathode ray tube or CRT is a specialized vacuum tube in which images are produced when an electron beam strikes a phosphorescent surface. Television sets, computers, automated teller machines, video game machines, video cameras, monitors, oscilloscopes and radar displays all contain cathode-ray tubes. Phosphor screens using multiple beams of electrons have allowed CRTs to display millions of colors. The first cathode ray tube scanning device was invented by the German scientist Karl Ferdinand Braun in 1897. Braun introduced a CRT with a fluorescent screen, known as the cathode ray oscilloscope. The screen would emit a visible light when struck by a beam of electrons. TV Tubes. TV tubes are basically cathode ray tubes. An electron beam is produced by an electrically-heated filament, and that beam is guided by two magnetic fields to a particular spot on the screen. The beam is moved so very quickly, that the eye can see not just one particular spot, but all the spots on the screen at once, forming a variable picture. The 2 magnetic fields are one for the vertical deflection, one for the horizontal deflection, and they are provided to the beam by external coils. Oscilloscope Tubes. Oscilloscope tubes are basically the same as TV tubes, but the beam is guided by two electrostatic fields provided by internal pieces of metal. It is a necessity because an oscilloscope uses a very large range of synchronisation frequencies for the deflection, when a TV set uses fixed frequencies, and it would be too difficult to drive large coils on a so large band of frequency. They are much deeper for the same screen size as a TV tube because the deflection angle is little. A TV tube has a deflection angle of 90° for the be

Electronics/DC Circuit Analysis. DC Circuit Analysis. In this chapter, capacitors and inductors will be introduced (without considering the effects of AC current.) The big thing to understand about Capacitors and Inductors in DC Circuits is that they have a transient (temporary) response. During the transient period, capacitors build up charge and stop the flow of current (eventually acting like infinite resistors.) Inductors build up energy in the form of magnetic fields, and become more conductive. In other words, in the steady-state (long term behavior), capacitors become open circuits and inductors become short circuits. Thus, for DC analysis, you can replace a capacitor with an empty space and an inductor with a wire. The only circuit components that remain are voltage sources, current sources, and resistors. Capacitors and Inductors at DC. DC steady-state (meaning the circuit has been in the same state for a long time), we've seen that capacitors act like open circuits and inductors act like shorts. The above figures show the process of replacing these circuit devices with their DC equivalents. In this case, all that remains is a voltage source and a lone resistor. (An AC analysis of this circuit can be found in the AC section.) Resistors. If a circuits contains only resistors possibly in a combination of parallel and series connections then an equivalent resistance is determined. Then Ohm's Law is used to determine the current flowing in the main circuit. A combination of voltage and current divider rules are then used to solve for other required currents and voltages. Simplify the following: The circuit in Figure 1 (a) is very simple if we are given R and V, the voltage of the source, then we use Ohm's Law to solve for the current. In Figure 1 (b) if we are given R1, R2 and V then we combine the resistor into an equivalent resistors noting that are in series. Then we solve for the current as before using Ohm's Law. In Figure 1 (c) if the resistors are labeled clockwise from the top resistor R1, R2 and R3 and the voltage source has the value: V. The analysis proceeds as follows. This is the formula for calculating the equivalent resistance of series resistor. The current is now calculated using Ohm's Law. If the voltage is required across the third resistor then we can use voltage divider rule. Or alternatively one could use Ohm's Law together with the current just calculated. In Figure 2 (a) if the Resistor nearest the voltage source is R1 and the other resistor R2. If we need to solve for the current i. then we proceed as before. First we calculate the equivalent resistance then use Ohm's Law to solve for the current. The resistance of a parallel combination is: So the current, "i", flowing in the circuit is, by Ohm's Law: If we need to solve for current through R2 then we can use current divider rule. But it would probably have been simpler to have used the fact that V most be dropped across R2. This means that we can simply use Ohm's Law to calculate the current through R2. The equation is just equation 1. In Figure 2 (b) we do exactly the same thing except this time there are three resistors this means that the equivalent resistance will be: Using this fact we do exactly the same thing. (a) (b) (c) In Figure 3 (a), if the three resistors in the outer loop of the circuit are R1, R2 and R3 and the other resistor is R4. It is simpler to see what is going on if we combine R2 and R3 into their series equivalent resistance formula_10. It is clear now that the equivalent resistance is R1 in series with the parallel combination of formula_10 and R4. If we want to calculate the voltage across the parallel combination of R4 and formula_12 then we just use voltage divider. If we want to calculate the current through R2 and R3 then we can use the voltage across formula_14 and Ohm's law. Or we could calculate the current in the main circuit and then use current divider rule to get the current. In Figure 3 (b) we take the same approach simplifying parallel combinations and series combinations of resistors until we get the equivalent resistance. In Figure 3 (c) this process doesn't work then because there are resistors connected in a delta this means that there is no way to simplify this beyond transforming them to a star or wye connection. Note: To calculate the current draw from the source the equivalent resistance always must be calculated. But if we just need the voltage across a series resistor this may be necessary. If we want to calculate the current in parallel combination then we must use either current divider rule or calculate the voltage across the resistor and then use Ohm's law to get the current. The second method will often require less work since the current flowing from the source is required for the use of current divider rule. The use of current divider rule is much simpler in the case when the source is a current source because the value of the current is set by the current source. The above image shows three points 1, 2, and 3 connected with resistors "R1", "R2", and "R3" with a common point. Such a configuration is called a "star network" or a "Y-connection". The above image shows three points 1, 2, and 3 connected with resistor "R12", "R23", and "R31". The configuration is called a "delta network" or "delta connection". We have seen that the series and parallel networks can be reduced by the use of simple equations. Now we will derive similar relations to convert a star network to delta and vice versa. Consider the points 1 and 2. The resistance between them in the star case is simply "R1 + R2" For the delta case, we have "R12 || (R31 + R23) " We have similar relations for the points 2, 3 and 3, 1. Making the substitution "r1= R23" etc., we have, simplifying, formula_16 formula_17 in the most general case. If all the resistances are equal, then "R = r/3".

Bicycles/Tire patches. Patches of rubber for repairing inner tubes. They are available in three varieties: regular, feather edged and self adhesive. All are available in a variety of sizes. Regular. These are simply cut or stamped from a sheet of material. The side to be glued is usually slightly sticky, and protected by a plastic film which should be removed just before use. As they are the same thickness all over there will be a sudden change in stiffness of the repaired inner tube at the edge of the patch. This can cause the patch to start peeling from the inner tube if not well stuck. Feather Edge. These are as thick as regular patches in the middle, but get thinner towards the edge. This makes it less likely that the patch will peel from the inner tube if not well stuck. They usually have a metal foil protecting the side to be glued, and a plastic film to support the thin edges while the patch is placed and pressed down. Self Adhesive. These are patches of either of the above types, but do not need a separate glue. Simply roughen the tube, peel the protective film from the patch and stick.

Perl Programming/Functions. A Perl function is a grouping of code that allows it to be easily used repeatedly. Functions are a key organizing component in all but the smallest of programs. Introduction. So far we have written just a few lines of Perl at a time. Our example programs have started at the top of a file and proceeded to the bottom of the file, with a little jumping around using control-flow keywords like if, else and while. In many instances, though, it is useful to add another layer of organization on our programs. For example, when something goes wrong in a typical program, it prints an "error message". The code to print an error message might look like print STDOUT "something went wrong!\n"; Other messages might look slightly different, like print STDOUT "something ELSE went wrong!\n"; We could pepper our code with hundreds of lines just like this, and things would work just fine…for a while. But sooner or later, we'd want to separate out error messages from "status messages" that report harmless information about our program. To do this, we could prefix all of our error messages with the word "ERROR" and all of the status messages with "STATUS". The code for a typical error messages would change to print STDOUT "ERROR: something went wrong!\n"; The problem is, with hundreds of error messages, it would be a hassle to go change them all. This is where subroutines can help out. The Wikipedia defines a subroutine, as "a sequence of instructions that perform a specific task, as part of a larger program. Subroutines can be called from different locations in a program, thus allowing programs to access the subroutine repeatedly without the subroutine's code having been written more than once." All that gobblydeegook means we can wrap up our error message code in one place, as follows: sub print_error_message { my($message) = @_; print STDOUT "ERROR: " . $message . "\n"; Whenever something goes wrong in our program, we can activate, or "call", this subroutine, with whatever message we prefer: print_error_message("something bad happened"); print_error_message("something really horrible happened"); print_error_message("something sort of annoying happened"); And see messages such as ERROR: something bad happened If we need to change the formatting of our error message, say, to include a few exclamation points, it's simple enough to change the subroutine: sub print_error_message { my($message) = @_; print STDOUT "ERROR: " . $message . "!!!\n"; This has been an admittedly simple example, and subroutines have a few other advantages, but that's it in a nutshell. Type it in one place, fix it in one place, change it in one place. Now for a little more detail about subroutines. Much of what follows will use the following subroutine, which, if it isn't obvious already, adds two numbers together, and returns their sum. sub add_two_numbers { my($x, $y) = @_; my $sum = $x + $y; return $sum; Parts of a subroutine. Name. sub add_two_numbers { my($x, $y) = @_; my $sum = $x + $y; return $sum; The first line of a function begins with the keyword sub followed by the function's name. Any string of letters and numbers that is not a reserved Perl word (such as for, while, if and else) is a valid function name. Subroutines with names that describe what the subroutine does make for easier-to-read programs. Prototype (rarely used). sub add_two_numbers($$) { The optional ($$) specifies how many arguments this subroutine expects. ($$) says "this function requires two scalar values". Perl prototypes are NOT what most people with experience with other languages expect them to be: instead, the prototypes alter the context of the parameters to the subroutine. It is possible to disable a prototype by calling the function with a leading &, but this is not recommended. Prototypes are used only seldomly, as it's much easier to use the normal method of passing parameters. Body. The body of a subroutine does the "work", and consists of three primary sections. Reading arguments. The pieces of information handed to a subroutine are called "arguments" or "actual parameters". For instance, in add_two_numbers(3, 4); 3 and 4 are the arguments to the subroutine add_two_numbers. Perl passes arguments to a subroutine in an array represented by @_. Usually it is more convenient to give meaningful names to these arguments, so the first line of a function often looks like sub add_two_numbers { my($x, $y) = @_; # reading parameters my $sum = $x + $y; return $sum; that puts the contents of @_ into two variables named $x and $y. $x and $y are called "formal parameters". The distinction between formal parameters (arguments) and actual parameters is subtle and for the most part unimportant. It is described somewhat in the wikipedia article Parameter (computer science). Be careful not to confuse the special variable $_ and @_, the array of arguments passed to a function. Some subroutines don't require any arguments, for example sub hello_world { print STDOUT "Hello World!\n"; will print "Hello World" to STDOUT. This subroutine doesn't need any extra information about how to do its job and accordingly doesn't need any arguments. Most modern programming languages save programmers the trouble of explicitly breaking the argument array into variables. Unfortunately, Perl does not. Which, on the other hand, makes writing subroutines with variable number of arguments very easy. In programming context "parameter" means almost the same thing as argument (see parameter for details). The two are often confused with no loss of understanding. Important note: global and local variables. Unlike programming languages such as C or Java, all variables created or used within Perl subroutines are, by default, "global" variables. This means that any piece of your program outside of your subroutine may modify these variables, and that your subroutine may be, unknowingly, modifying variables that it has no business modifying. In small programs this is often convenient, but as programs get longer this often leads to complexity and is considered poor practice. The best way to avoid this trap is to place the keyword my in front of all your variables the first time they appear. This tells Perl that you only want these variables to be available inside the nearest enclosing group of curly braces. In effect, these "local" variables act as a "scratch space" for use within your subroutines that disappears when the subroutine returns. The line use strict; at the top of your program will instruct Perl to "force" you to use my in front of your variables, to prevent you from accidentally creating global variables. An alternative to my that you may see in some older Perl programs is the local keyword. local is somewhat similar to my, but is more complicated to deal with. It's better to stick with my in your own programs. "Scope" describes whether a variable is local or global, and a couple of other complexities. See scope for a technical discussion. The interesting part of a subroutine. In the very middle of the subroutine you're likely to find the more interesting "guts" of it all. In our add_two_numbers subroutine, this is the part that actually does the adding sub add_two_numbers { my($x, $y) = @_; my $sum = $x + $y; # the interesting part return $sum; In this middle section, you can do just about anything your heart desires, arithmetic, printing to files, or even calling other subroutines. The return statement. Finally, some subroutines "return" some piece of information, called the "return value", using the return keyword. sub add_two_numbers { my($x, $y) = @_; my $sum = $x + $y; return $sum; # the return statement For example $sum = add_two_numbers(4, 5); will set $sum to 9 (the sum of 4 and 5). Invoking subroutines. Subroutines may be declared anywhere within a Perl program. They may be invoked like so: add_two_numbers(4, 5); # the safest approach add_two_numbers 4, 5; # only if predeclared &add_two_numbers(4, 5); # older Perl syntax, but still valid If no "&" prefix is used, parentheses are required unless the subroutine has been predeclared. Functions calling functions. On their own, functions provide a major stepping stone towards good code, but combining functions together really unleashes their power. As you might expect, calling a function from "within" another function doesn't look any different from calling a function from the part of your program that is sitting outside any curly-braces. This function adds two numbers, then multiplies them by 3. Bear with us on the uselessness of these functions. As you build your own programs to solve your unique problems, you'll see their usefulness immediately. sub add_two_numbers_and_mult_by_three { my($x, $y) = @_; # read parameters my $sum = add_two_numbers($x, $y); # add x and y, put result in sum my $sum_times_three = $sum*3; # multiply by three return $sum_times_three; # return result The line my $sum = add_two_numbers($x, $y); # add x and y, put result in sum calls our function add_two_numbers and puts the result into our $sum variable. Easy stuff, huh? In this function, we've actually written much more code than we need. It could be pared down to something much smaller, but equally readable: sub add_two_numbers_and_mult_by_three { my($x, $y) = @_; # read parameters return 3*add_two_numbers($x, $y); # add x and y, mult by 3, return Functions calling themselves — recursion. We've seen functions calling "other" functions, but one neat concept in programming is when functions call "themselves". This is called recursion. At first it seems like this might cause a so-called infinite loop, but it's really quite standard programming. In math, the factorial function, multiplies a positive integer by each positive integer less than itself. For example, "5 factorial" (usually written 5!) is calculated by multiplying 5 times 4 times 3 times 2 times 1. Of course, the 1 doesn't change the result. The factorial function is useful in calculating things like the number of different possible ways to seat your relatives at the dinner table. Factorial makes a natural example for recursion, though it can be written just as easily with a while loop. sub factorial { my($num) = @_; if ($num == 1) { return 1; # stop at 1, factorial doesn't multiply times zero } else { return $num*factorial($num - 1); # call factorial function recursively The self-referential line here is return $num*factorial($num - 1); # call factorial function recursively which calls the factorial function from "within" the factorial function. This would go on forever, but we have a sort of stop-sign that prevents it: if ($num == 1) { return 1; # stop at 1, factorial doesn't multiply times zero This stops the sequence of calls to factorial and prevents the never-ending infinite loop. Written in a sort of longhand factorial(5) = 5*factorial(4) = 5*4*factorial(3) = 5*4*3*factorial(2) = 5*4*3*2*factorial(1) = 5*4*3*2*1 = 120 We've just barely touched on recursion. For some programming problems, it is a very natural solution. For others, it's a little… unnatural. Suffice it to say, it's a tool every programmer should carry in his or her belt. Functions vs procedures vs subroutines. While reading programming literature and talking to programmers, you might run across the three terms "functions", "procedures" and "subroutines". Most of the time these are used interchangeably, but for the purists out there: A function always returns a value, and, given the same arguments, always returns the same value. Functions are a lot like the functions you may have used in math class. A procedure, unlike a function, may return no value at all. Unlike a function, a procedure often interacts with its external environment beyond its argument array. For instance, a procedure may read and write to files. A subroutine refers to a sequence of instructions that may be either a function or a procedure. The official grammar. The syntax for defining a subroutine is: sub NAME PROTOTYPE ATTRIBUTES BLOCK If this makes any sense to you, you probably don't need to be reading this book ;) </noinclude>

Organic Chemistry/Introduction to reactions/Diels-Alder reaction. The Diels-Alder reaction, named after the German chemists who developed it, is a method for producing simple ring compounds. <br>"Mechanism of a reaction between a diene and a dienophile" In the Diels-Alder reaction, a conjugated diene reacts with a dienophile. The dienophile is named for its affinity to react with the diene. Because dienes and dienophiles are often gases, this reaction usually takes place in a sealed container at elevated pressure and temperature. The temperature required by this reaction can be reduced by the presence of electron-witdrawing groups attached to the dienophile. In some cases, such as furan and maleic anhydride, or cyclopentadiene and acrolein, the reaction will take place at room temperature in an ethoxyethane solution. The reaction shown above is highly unlikely due to the lack of substituents on the diene and dienophile. The dienophile will usually have a carbonyl, nitro, or other electron-withdrawing group attached. Any substituents attached to the diene or dienophile will end up in the final product. <br>"Notice that the product retains the same basic cyclohexene structure. The substituents simply extend from the ring." Alkynes can also react as dienophiles. Diels-Alder reactions can sometimes reverse themselves through Retro-Diels-Alder reactions. For example, dicyclopentadiene can be cracked to form 1,3-cyclopentadiene by thermal dissociation. Retro reactions occur under situations where the fragments are stable by themselves.

XML - Managing Data Exchange/Parsing XML files. In the earlier chapters we were taught how to create XML files in detail. This involved the development of XML documents, Style sheets and Schema and their validation. In this chapter, we will focus on different approaches for parsing XML files and when to use them. But first, it is time to refresh what we have learned about parsing. The Process of Parsing XML files. One goal of the XML format was to enhance raw data formats like plain text by including detailed descriptions of the meaning of the content. Now, in order to be able to read XML files, we use a parser which basically exposes the document’s content through a so-called API (application programming interface). In other words, a client application accesses the content of the XML document through an interface, instead of having to interpret the XML code on its own! Simple Text Parsing. One way to extract data from an XML document is simple text parsing – browsing all characters in the document and check for a desired pattern: <house> <value><int>150,000</int></value> </house> Let’s say we are interested in the value of the house. Using straight text parsing, we would scan the file for the character sequence and call it the start pattern. Then, we would further scan the document for the end pattern (i.e. ). Finally, we declare the text string in between these two patterns to be the value of the surrounding tag. Why it doesn't work that way. Obviously, this approach is not suitable for extracting information from large and complex XML documents, since we would have to know exactly what the file looks like and where the information needed is located. From a more general point of view, the structure and semantics of an XML file is determined by the makeup of the document, its tags and attributes – hence, we need a device that is able to recognize and understand this structure and can point out any errors in it. Moreover, it has to provide the content of the document through an interface, so that other applications can access it without difficulty. This device is known as an XML parser. What a parser does. Almost all programs that need to process XML documents use an XML parser to extract the information stored in the XML document in order to avoid any of the difficulties that occur when reading and interpreting raw XML data. The parser usually is a class library (e.g. a set of Java class files) that reads a given document and checks if it is well-formed according to the W3C specification. Then, any client software can use methods of the interface provided by the parser API to access the information the parser retrieved from the XML file. All in all, the parser shields the user from dealing with the complex details of XML like assembling information distributed over several XML files, checking for well-formedness constraints, and so on. Parsing: an Example. To illustrate more clearly what parsing an XML file really means, the following example was created which contains information about some cities. It also keeps track of who is on vacation and demonstrates the parsing process with the currently most common parsing methods. Example: "cities.xml". <?xml version="1.0" encoding="UTF-8" ?> <cities> <city vacation="Sam"> <cityName>Atlanta</cityName> <cityCountry>USA</cityCountry> </city> <city vacation="David"> <cityName>Sydney</cityName> <cityCountry>Australia</cityCountry> </city> <city vacation="Pune"> <cityName>Athens</cityName> <cityCountry>Greece</cityCountry> </city> </cities> Based on the information stored in this XML document, we can easily check who is on vacation and where. The parser will read the file using one of the various techniques presented later in this chapter. This process is very complicated and prone to errors of all kinds. Luckily, we will never have to write code for it, because there are plenty of free, fully-functional parsers on the Web. All we do is download a parser class library and access the XML document through the interface provided by the parser software. With more recent builds of Java, most parsers do not even have to be downloaded. In other words, we use the functions or methods included in the class library for extracting the information. Basically, a parser reads the XML document and tries to recognize the structure of the file itself while checking for errors. It simply checks for start/end tags, attributes, namespaces, prefixes, and so on. Then, the client software can access the information derived from this structure using methods provided by the parser software (i.e. the interface). The best way to learn about the functionality of a parser is to actually use them; therefore, the next section demonstrates the different methods of parsing. Parser APIs (Application Programming Interface). Overview. There are two “traditional” approaches that dominate the market right now, an event-based push-model as represented by SAX (Simple API for XML) and a tree-based model using the DOM (document object model) approach. However, there is a movement towards newer approaches and techniques that try to overcome the flaws inherent in these traditional models – an event-based pull-model and a “cursor model”, such as VTD-XML, which allows us to browse the XML document just like in the tree-based approach, but simpler and easier to use. SAX (Simple API for XML). Description. The push model, typically the exemplified by SAX (www.saxproject.org) is the “gold standard” of XML parsing, since it is probably the most complete and accurate method so far. The SAX classes provide an interface between the input streams from which XML documents are read and the client software which receives the data made available by the parser. The parser browses through the whole document and fires events every time it recognizes an XML construct (e.g. it recognizes a start tag and fires an event – the client software is notified and can use this information… or not). Evaluation. The advantage of such a model is that we don’t need to store the whole XML document in memory, since we are only reading one piece of information at a time. If you recall that the XML structure is a set of nodes of various types (like an element node) – parsing the document with a SAX parser means going through each node one at a time. This makes it possible to read even very large XML documents in a memory-efficient way. However, the fact that the parser only provides information about the node currently read also implies that the programmer of the client software is in charge of saving certain information in a separate data structure (e.g. the parents or children of the currently processed node). Moreover, the SAX approach is pretty much read-only, since it is hard to modify the XML structure when we do not have some sort of global view. In fact, the parser is in control of what is read when. The user can only wait until a certain event has occurred and then use the information stored in the currently processed node. Example: "TGSAXParser.java". As mentioned before, the best way to fully understand the concept of the parsing process is to actually use it. In the following code sample, the information about the name and country of the cities that people are vacationing in will be displayed. The SAX API that is part of the Xerces parser package was used for the implementation ((Xerces 2 Homepage): // import the basic SAX API classes import org.xml.sax.*; import org.xml.sax.helpers.*; import java.io.*; public class TGSAXParser extends DefaultHandler public boolean onVacation = false; // what to do when a start-element event was triggered public void startElement(String uri, String name, String qName, Attributes atts) // stores the string in the XML file String vacationer = atts.getValue("vacation"); String cityName = atts.getValue("cityName"); String cityCountry = atts.getValue("cityCountry"); // if the start tag is "city" set vacationer to true if (qName.equals("city") && (vacationer != null)) onVacation = true; System.out.print("\n" + vacationer + " is on vacation in "); if (qName.equals("cityName") && onVacation) if (qName.equals("cityCountry") && onVacation) /**This method is used to stop printing information once the element has *been read. It will also reset the onVacation variable for the next *element. public void endElement(String uri, String name, String qName) //reset flag if (qName.equals("city")) onVacation = false; /**This method is triggered to store and print the values between *the XML tags. It will only print those values if onVacation == true. public void characters(char[] ch, int start, int length) if (onVacation) for (int i = start; i < start + length; i++) System.out.print(ch[i]); public static void main(String[] args) System.out.println("People on vacation in the following cities:"); try // create a SAX parser from the Xerces package XMLReader xml = XMLReaderFactory.createXMLReader(); TGSAXParser handler = new TGSAXParser(); xml.setContentHandler(handler); xml.setErrorHandler(handler); FileReader r = new FileReader("cities.xml"); xml.parse(new InputSource(r)); catch (SAXException se) System.out.println("XML Parsing Error: " + se); catch (IOException io) System.out.println("File I/O Error: " + io); The codice_1: As mentioned before, SAX is completely event-driven. Therefore, we need a handler that “listens” to the input stream coming from the input file ("cities.xml" in this case). The SAX API provides interface classes, which we have to extend with our own code to read our own specific XML document. In order to include our code in the SAX API, we just have to extend the DefaultHandler interface with our own class and set the content handler to our custom handler class (which consists of three methods: startElement, endElement and characters) The codice_2 and codice_3 methods: These methods are invoked whenever the SAX parser finds a start or end tag respectively. The SAX API provides blank stubs for both methods and we have to fill them with code of our own. In this case, we want our program to do something whenever the codice_4 attribute is set, so we set a Boolean variable to true whenever we find such an element and process the node by printing out the character sequence in between the start and end tag. The codice_5 method is automatically called whenever a codice_6 and codice_7 event was triggered, but prints out the character string only if the codice_8 attribute is set. DOM (Document Object Model). Description. The other popular approach is the tree-based model as represented by the DOM (document object model, see W3C Recommendation). This method actually works similarly to a SAX parser, since it reads the XML document from an input stream by browsing through the file and recognizing XML structures. This time, instead of returning the content of the document in a series of small fragments, the DOM method maps the XML hierarchy to a DOM tree object that contains everything from the original XML document. Everything from elements, comments, textual information or processing instructions is stored in the tree object as nodes, starting with the document itself as the root node. Now that all the information we need is stored in memory, we access the data by using methods provided by the parser software to read or modify objects within the tree. This facilitates random access to the content of the XML document and provides the possibility to modify the data it contains or even create new XML files by transforming a DOM back to an XML document. Evaluation. However, the major downside of this approach is that it requires much more memory and is therefore not suitable for situations where large XML files are used. More importantly, it is somewhat more complex than the simplistic SAX method even for small and simple problems. Example: "MyDOMParser.java". In the following code sample, a list of cities with people on vacation is again created but this time with the tree-based approach: // import all necessary DOM API classes import org.apache.xerces.parsers.*; import org.apache.xerces.dom.*; import org.w3c.dom.*; public class MyDOMParser{ public static void main(String[] args) { System.out.println("People on vacation in the following cities:"); try { // creates a DOM parser object DOMParser parser = new DOMParser(); parser.parse("cities.xml"); // stores the tree object in a variable org.w3c.dom.Document doc = parser.getDocument(); // returns a list of all city elements in my city list NodeList list = doc.getElementsByTagName("city"); // now, for every element in the city list, check if the // "vacation" attribute is set and if yes, print out the // information about the vacationer. for(int i = 0, length = list.getLength(); i < length; i++){ Element city = (Element)list.item(i); Attr vacationer = city.getAttributeNode("vacation"); if(vacationer!= null){ String v = vacationer.getValue(); System.out.print(v + " is vacationing in "); // grab information about city name and country // directly from the DOM tree object ParentNode cityname = (ParentNode) doc.getElementsByTagName("cityName").item(0); ParentNode country = (ParentNode) doc.getElementsByTagName("cityCountry").item(0); System.out.println(cityname.getTextContent() + ", " + country.getTextContent()); } catch (Exception e) { System.out.println(e.getMessage()); codice_9: Once we parsed the XML document, the tree object is temporarily stored in the parser variable. In order to work with the DOM object, we have to create a variable holding it (of type codice_10). Then, we create a list of nodes holding all elements with the tag name . The parser finds these nodes by browsing through the DOM tree. Then, we just go through each one of the city-elements and check if the vacation attribute is set and display all the information about the vacationer if so. Xerces provides a helpful method called codice_11 that lets us directly access the text node of an element node, avoiding all difficulties emerging from unneeded white space and the like. Summary. Choosing an API at the beginning of your XML project is a very important decision. Once you decide which one to use, it is easy to try different vendors without having much trouble, but switching to a different API will be a very time-consuming and costly process, since you will have to redesign your whole program code. The SAX API is a widely accepted and well-working parser that is easy to implement and works especially well with streaming content (e.g. an online XML source). Because it is a read-only API, you would not be able to modify the underlying XML data source. Since it only reads one node at a time, it is very memory-efficient and fast. However, this implies that your application expects the information to be close together and ordered. If you want to randomly access the entire document at any point of time, then the DOM approach might be a better choice for you. The DOM API is more complex and harder to implement, but gives you full control over the whole document and lets you modify the data, also. However, it reads the whole XML document into memory, so the DOM API is not suitable for projects with very large XML files. Exercise. Recommended optional exercise. "Use the code sample for the SAX and DOM parser from this chapter and play around with it. You probably want to print out different nodes or add more constraints. This absolutely optional, but will give you an idea of the main differences between SAX and DOM." Now for the exercise. "TO HELP YOU" download this, it provides a structure of the problem so that you can more easily run the app in NetBeans 5.0. If you’re interested in using Xerces – just download the following file: http://www.apache.org/dist/xml/xerces-j/Xerces-J-bin.2.8.0.zip If the above link is dead. Go to http://www.apache.org/dist/xml/xerces-j/ and download the latest zip binary file. It should be in the format of "Xerces-J-bin.#.#.#.zip" Then put the content into the \lib\ext subfolder of your NetBeans directory and start up NetBeans IDE. Now, the Xerces package is successfully installed on your machine.

Calculus/Partial differential equations. First order. Any partial differential equation of the form where "h"1, "h"2 … "h"n, and "b" are all functions of both "u" and Rn can be reduced to a set of ordinary differential equations. To see how to do this, we will first consider some simpler problems. Special cases. We will start with the simple PDE Because "u" is only differentiated with respect to "z", for any fixed "x" and "y" we can treat this like the ODE, "du/dz=u". The solution of that ODE is "cez", where "c" is the value of "u" when "z=0", "for the fixed x and y" Therefore, the solution of the PDE is Instead of just having a constant of integration, we have an arbitrary function. This will be true for any PDE. Notice the shape of the solution, an arbitrary function of points in the "xy", plane, which is normal to the 'z' axis, and the solution of an ODE in the 'z' direction. Now consider the slightly more complex PDE where h can be any function, and each "a" is a real constant. We recognize the left hand side as being a·∇, so this equation says that the differential of "u" in the a direction is "h(u)". Comparing this with the first equation suggests that the solution can be written as an arbitrary function on the plane normal to a combined with the solution of an ODE. Remembering from Calculus/Vectors that any vector r can be split up into components parallel and perpendicular to "a", we will use this to split the components of r in a way suggested by the analogy with (1). Let's write and substitute this into (2), using the chain rule. Because we are only differentiating in the a direction, adding any function of the perpendicular vector to "s" will make no difference. First we calculate grad "s", for use in the chain rule, On making the substitution into (2), we get, which is an ordinary differential equation with the solution The constant "c" can depend on the perpendicular components, but not upon the parallel coordinate. Replacing "s" with a monotonic scalar function of "s" multiplies the ODE by a function of "s", which doesn't affect the solution. "Example:" For this equation, a is (1, -1), "s=x-t", and the perpendicular vector is (x+t)(1, 1). The reduced ODE is "du/ds=0" so the solution is To find "f" we need initial conditions on "u". Are there any constraints on what initial conditions are suitable? Consider, if we are given Similarly, if we are given "u" on any curve which the lines "x"+"t"="c" intersect only once, and to which they are not tangent, we can deduce "f". For any first order PDE with constant coefficients, the same will be true. We will have a set of lines, parallel to "r"=at, along which the solution is gained by integrating an ODE with initial conditions specified on some surface to which the lines aren't tangent. If we look at how this works, we'll see we haven't actually used the constancy of a, so let's drop that assumption and look for a similar solution. The important point was that the solution was of the form "u"="f"("x"("s"),"y"("s")), where ("x"("s"),"y"("s")) is the curve we integrated along -- a straight line in the previous case. We can add constant functions of integration to "s" without changing this form. Consider a PDE, For the suggested solution, "u"="f"("x"("s"),"y"("s")), the chain rule gives Comparing coefficients then gives so we've reduced our original PDE to a set of simultaneous ODE's. This procedure can be reversed. The curves ("x"("s"),"y"("s")) are called "characteristics" of the equation. "Example:" Solve formula_14 given "u"="f"("x") for "x"≥0 The ODE's are subject to the initial conditions at "s"=0, This ODE is easily solved, giving so the characteristics are concentric circles round the origin, and in polar coordinates "u"("r","θ")="f"("r") Considering the logic of this method, we see that the independence of "a" and "b" from "u" has not been used either, so that assumption too can be dropped, giving the general method for equations of this "quasilinear" form. Quasilinear. Summarising the conclusions of the last section, to solve a PDE subject to the initial condition that on the surface, ("x"1("r"1,…,"r""n"-1, …"x"n("r"1,…,"r""n"-1), u=f("r"1,…,"r""n"-1) --this being an arbitrary paremetrisation of the initial surface-- Both the second and third steps may be troublesome. The set of ODEs is generally non-linear and without analytical solution. It may even be easier to work with the PDE than with the ODEs. In the third step, the "r""i" together with "s" form a coordinate system adapted for the PDE. We can only make the inversion at all if the Jacobian of the transformation to Cartesian coordinates is not zero, This is equivalent to saying that the vector ("a""1", &hellip:, "a""n") is never in the tangent plane to a surface of constant "s". If this condition is not false when "s"=0 it may become so as the equations are integrated. We will soon consider ways of dealing with the problems this can cause. Even when it is technically possible to invert the algebraic equations it is obviously inconvenient to do so. Example. To see how this works in practice, we will <br> a/ consider the PDE, with generic initial condition, Naming variables for future convenience, the corresponding ODE's are subject to the initial conditions at τ=0 These ODE's are easily solved to give These are the parametric equations of a set of straight lines, the characteristics. The determinant of the Jacobian of this coordinate transformation is This determinant is 1 when "t"=0, but if "fr" is anywhere negative this determinant will eventually be zero, and this solution fails. In this case, the failure is because the surface formula_28 is an envelope of the characteristics. For arbitrary "f" we can invert the transformation and obtain an implicit expression for "u" If "f" is given this can be solved for "u". 1/ formula_30, The implicit solution is This is a line in the "u-x" plane, rotating clockwise as "t" increases. If "a" is negative, this line eventually become vertical. If "a" is positive, this line tends towards "u"=0, and the solution is valid for all "t". "For f"("x","y")="x"2 The implicit solution is: This solution clearly fails when formula_33, which is just when formula_28 . For any "t >" 0 this happens somewhere. As "t" increases this point of failure moves toward the origin. Notice that the point where "u"=0 stays fixed. This is true for any solution of this equation, for all values of "f" We will see later that we can find a solution after this time, if we consider discontinuous solutions. We can think of this as a shockwave. 3/ formula_35<br>The implicit solution is and we can not solve this explitely for "u". The best we can manage is a numerical solution of this equation. We can also consider the closely related PDE The corresponding ODE's are subject to the initial conditions at τ=0 These ODE's are easily solved to give Writing "f" in terms of "u", "s", and τ, then substituting into the equation for "x" gives an implicit solution It is possible to solve this for "u" in some special cases, but in general we can only solve this equation numerically. However, we can learn much about the global properties of the solution from further analysis Characteristic initial value problems. What if initial conditions are given on a characteristic, on an envelope of characteristics, on a surface with characteristic tangents at isolated points? Discontinuous solutions. So far, we've only considered smooth solutions of the PDE, but this is too restrictive. We may encounter initial conditions which aren't smooth, e.g. If we were to simply use the general solution of this equation for smooth initial conditions, we would get which appears to be a solution to the original equation. However, since the partial differentials are undefined on the characteristic "x"+"ct"=0, so it becomes unclear what it means to say that the equation is true at that point. We need to investigate further, starting by considering the possible types of discontinuities. If we look at the derivations above, we see we've never use any second or higher order derivatives so it doesn't matter if they aren't continuous, the results above will still apply. The next simplest case is when the function is continuous, but the first derivative is not, e.g. |"x"|. We'll initially restrict ourselves to the two-dimensional case, "u"("x", "t") for the generic equation. Typically, the discontinuity is not confined to a single point, but is shared by all points on some curve, ("x"0("s"), "t"0("s") ) Then we have We can then compare "u" and its derivatives on both sides of this curve. It will prove useful to name the "jumps" across the discontinuity. We say Now, since the equation (1) is true on both sides of the discontinuity, we can see that both "u"+ and "u"-, being the limits of solutions, must themselves satisfy the equation. That is, Subtracting then gives us an equation for the jumps in the differentials We are considering the case where "u" itself is continuous so we know that ["u"]=0. Differentiating this with respect to "s" will give us a second equation in the differential jumps. The last two equations can only be both true if one is a multiple of the other, but multiplying "s" by a constant also multiplies the second equation by that same constant while leaving the curve of discontinuity unchanged, hence we can without loss of generality define "s" to be such that But these are the equations for a characteristic, i.e. discontinuities propagate along characteristics. We could use this property as an alternative definition of characteristics. We can deal similarly with discontinuous functions by first writing the equation in "conservation form", so called because conservation laws can always be written this way. Notice that the left hand side can be regarded as the divergence of ("au", "bu"). Writing the equation this way allows us to use the theorems of vector calculus. Consider a narrow strip with sides parallel to the discontinuity and width "h" We can integrate both sides of (1) over R, giving Next we use Green's theorem to convert the left hand side into a line integral. Now we let the width of the strip fall to zero. The right hand side also tends to zero but the left hand side reduces to the difference between two integrals along the part of the boundary of "R" parallel to the curve. The integrals along the opposite sides of "R" have different signs because they are in opposite directions. For the last equation to always be true, the integrand must always be zero, i.e. Since, by assumption ["u"] isn't zero, the other factor must be, which immediately implies the curve of discontinuity is a characteristic. Once again, "discontinuities propagate along characteristics." Above, we only considered functions of two variables, but it is straightforward to extend this to functions of "n" variables. The initial condition is given on an "n"-1 dimensional surface, which evolves along the characteristics. Typical discontinuities in the initial condition will lie on a "n"-2 dimensional surface embedded within the initial surface. This surface of discontinuity will propagate along the characteristics that pass through the initial discontinuity. The jumps themselves obey ordinary differential equations, much as "u" itself does on a characteristic. In the two dimensional case, for "u" continuous but not smooth, a little algebra shows that while "u" obeys the same equation as before, We can integrate these equations to see how the discontinuity evolves as we move along the characteristic. We may find that, for some future "s", ["u""x"] passes through zero. At such points, the discontinuity has vanished, and we can treat the function as smooth at that characteristic from then on. Conversely, we can expect that smooth functions may, under the right circumstances, become discontinuous. To see how all this works in practice we'll consider the solutions of the equation for three different initial conditions. The general solution, using the techniques outlined earlier, is "u" is constant on the characteristics, which are straight lines with slope dependent on "u". First consider "f" such that While "u" is continuous its derivative is discontinuous at "x"=0, where "u"=0, and at "x"=a, where "u"=1. The characteristics through these points divide the solution into three regions. All the characteristics to the right of the characteristic through "x"=a, "t"=0 intersect the "x"-axis to the right of "x"=1, where "u"=1 so "u" is 1 on all those characteristics, i.e whenever "x"-"t">a. Similarly the characteristic through the origin is the line "x"=0, to the left of which "u" remains zero. We could find the value of "u" at a point in between those two characteristics either by finding which intermediate characteristic it lies on and tracing it back to the initial line, or via the general solution. Either way, we get At larger "t" the solution "u" is more spread out than at "t"=0 but still the same shape. We can also consider what happens when "a" tends to 0, so that "u" itself is discontinuous at "x"=0. If we write the PDE in conservation form then use Green's theorem, as we did above for the linear case, we get ["u"²] is the difference of two squares, so if we take "s=t" we get In this case the discontinuity behaves as if the value of "u" on it were the average of the limiting values on either side. However, there is a caveat. Since the limiting value to the left is "u"- the discontinuity must lie on that characteristic, and similarly for "u"+; i.e "the jump discontinuity must be on an intersection of characteristics", at a point where "u" would otherwise be multivalued. For this PDE the characteristic can only intersect on the discontinuity if If this is not true the discontinuity can not propagate. Something else must happen. The limit "a"=0 is an example of a jump discontinuity for which this condition is false, so we can see what happens in such cases by studying it. Taking the limit of the solution derived above gives If we had taken the limit of any other sequence of initials conditions tending to the same limit we would have obtained a trivially equivalent result. Looking at the characteristics of this solution, we see that at the jump discontinuity characteristics on which "u" takes every value between 0 and 1 all intersect. At later times, there are two slope discontinuities, at "x"=0 and "x"="t", but no jump discontinuity. This behaviour is typical in such cases. The jump discontinuity becomes a pair of slope discontinuities between which the solution takes all appropriate values. Now, lets consider the same equation with the initial condition This has slope discontinuities at "x"=0 and "x"=a, dividing the solution into three regions. The boundaries between these regions are given by the characteristics through these initial points, namely the two lines These characteristics intersect at "t"=a, so the nature of the solution must change then. In between these two discontinuities, the characteristic through "x"="b" at "t"=0 is clearly All these characteristics intersect at the same point, ("x","t")=("a","a"). We can use these characteristics, or the general solution, to write "u" for "t"<"a" As "t" tends to "a", this becomes a step function. Since "u" is greater to the left than the right of the discontinuity, it meets the condition for propagation deduced above, so for "t">"a" "u" is a step function moving at the average speed of the two sides. This is the reverse of what we saw for the initial condition previously considered, two slope discontinuities merging into a step discontinuity rather than vice versa. Which actually happens depends entirely on the initial conditions. Indeed, examples could be given for which both processes happen. In the two examples above, we started with a discontinuity and investigated how it evolved. It is also possible for solutions which are initially smooth to become discontinuous. For example, we saw earlier for this particular PDE that the solution with the initial condition "u"="x"² breaks down when 2"xt"+1=0. At these points the solution becomes discontinuous. Typically, discontinuities in the solution of any partial differential equation, not merely ones of first order, arise when solutions break down in this way and propagate similarly, merging and splitting in the same fashion. Fully non-linear PDEs. It is possible to extend the approach of the previous sections to reduce any equation of the form to a set of ODE's, for "any" function, "F". We will not prove this here, but the corresponding ODE's are If "u" is given on a surface parameterized by "r"1…"r""n" then we have, as before, "n" initial conditions on the "n", "x""i" given by the parameterization and "one" initial condition on "u" itself, but, because we have an extra "n" ODEs for the "u""i"'s, we need an extra "n" initial conditions. These are, "n"-1 consistency conditions, which state that the "u""i"'s are the partial derivatives of "u" on the initial surface, and "one" initial condition stating that the PDE itself holds on the initial surface. These "n" initial conditions for the "u""i" will be a set of algebraic equations, which may have multiple solutions. Each solution will give a different solution of the PDE. Example. Consider The initial conditions at τ=0 are and the ODE's are Note that the partial derivatives are constant on the characteristics. This always happen when the PDE contains only partial derivatives, simplifying the procedure. These equations are readily solved to give On eliminating the parameters we get the solution, which can easily be checked. abc Second order. Suppose we are given a second order linear PDE to solve formula_83 The natural approach, after our experience with ordinary differential equations and with simple algebraic equations, is attempt a factorisation. Let's see how for this takes us. We would expect factoring the left hand of (1) to give us an equivalent equation of the form and we can immediately divide through by "a". This suggests that those particular combinations of first order derivatives will play a special role. Now, when studying first order PDE's we saw that such combinations were equivalent to the derivatives along characteristic curves. Effectively, we changed to a coordinate system defined by the characteristic curve and the initial curve. Here, we have two combinations of first order derivatives each of which may define a different characteristic curve. If so, the two sets of characteristics will define a natural coordinate system for the problem, much as in the first order case. In the new coordinates we will have with each of the factors having become a differentiation along its respective characteristic curve, and the left hand side will become simply "u""rs" giving us an equation of the form If "A", "B", and "C" all happen to be zero, the solution is obvious. If not, we can hope that the simpler form of the left hand side will enable us to make progress. However, before we can do all this, we must see if (1) can actually be factored. Multiplying out the factors gives formula_87 On comparing coefficients, and solving for the α's we see that they are the roots of Since we are discussing real functions, we are only interested in real roots, so the existence of the desired factorization will depend on the discriminant of this quadratic equation. It can be shown that, just as for first order PDEs, discontinuities propagate along characteristics. Since elliptic equations have no real characteristics, this implies that any discontinuities they may have will be restricted to isolated points; i.e., that the solution is almost everywhere smooth. This is not true for hyperbolic equations. Their behavior is largely controlled by the shape of their characteristic curves. These differences mean different methods are required to study the three types of second equation. Fortunately, changing variables as indicated by the factorisation above lets us reduce any second order PDE to one in which the coefficients of the second order terms are constant, which means it is sufficient to consider only three standard equations. We could also consider the cases where the right hand side of these equations is a given function, or proportional to "u" or to one of its first order derivatives, but all the essential properties of hyperbolic, parabolic, and elliptic equations are demonstrated by these three standard forms. While we've only demonstrated the reduction in two dimensions, a similar reduction applies in higher dimensions, leading to a similar classification. We get, as the reduced form of the second order terms, where each of the "a""i"s is equal to either 0, +1, or -1. If "all" the "a""i"s have the "same sign" the equation is "elliptic" If "any" of the "a""i"s are "zero" the equation is "parabolic" If "exactly one" of the "a""i"s has the "opposite sign" to the rest the equation is "hyperbolic" In 2 or 3 dimensions these are the only possibilities, but in 4 or more dimensions there is a fourth possibility: "at least two" of the "a""i"s are "positive", and "at least two" of the "a""i"s are "negative". Such equations are called "ultrahyperbolic". They are less commonly encountered than the other three types, so will not be studied here. When the coefficients are not constant, an equation can be hyperbolic in some regions of the "xy" plane, and elliptic in others. If so, different methods must be used for the solutions in the two regions. Elliptic. Standard form, Laplace's equation: Quote equation in spherical and cylindrical coordinates, and give full solution for cartesian and cylindrical coordinates. Note averaging property Comment on physical significance, rotation invariance of laplacian. Hyperbolic. Standard form, wave equation: Solution, any sum of functions of the form These are waves. Compare with solution from separating variables. Domain of dependence, etc. Parabolic. The canonical parabolic equation is the diffusion equation: Here, we will consider some simple solutions of the one-dimensional case. The properties of this equation are in many respects intermediate between those of hyperbolic and elliptic equation. As with hyperbolic equations but not elliptic, the solution is well behaved if the value is given on the initial surface "t"=0. However, the characteristic surfaces of this equation are the surfaces of constant "t", thus there is no way for discontinuities to propagate to positive t. Therefore, as with elliptic equations but not hyberbolic, the solutions are typically smooth, even when the initial conditions aren't. Furthermore, at a local maximum of "h", its Laplacian is negative, so "h" is decreasing with "t", while at local minima, where the Laplacian will be positive, "h" will increase with "t". Thus, initial variations in "h" will be smoothed out as "t" increases. In one dimension, we can learn more by integrating both sides, Provided that "hx" tends to zero for large "x", we can take the limit as "a" and "b" tend to infinity, deducing so the integral of "h" over all space is constant. This means this PDE can be thought of as describing some conserved quantity, initially concentrated but spreading out, or diffusing, over time. This last result can be extended to two or more dimensions, using the theorems of vector calculus. We can also differentiate any solution with respect to any coordinate to obtain another solution. E.g. if "h" is a solution then so "hx" also satisfies the diffusion equation. Similarity solution. Looking at this equation, we might notice that if we make the change of variables then the equation retains the same form. This suggests that the combination of variables "x"²/"t", which is unaffected by this variable change, may be significant. We therefore assume this equation to have a solution of the special form then and substituting into the diffusion equation eventually gives which is an ordinary differential equation. Integrating once gives Reverting to "h", we find This last integral can not be written in terms of elementary functions, but its values are well known. In particular the limiting values of "h" at infinity are taking the limit as "t" tends to zero gives We see that the initial discontinuity is immediately smoothed out. The solution at later times retains the same shape, but is more stretched out. The derivative of this solution with respect to "x" is itself a solution, with "h" spreading out from its initial peak, and plays a significant role in the further analysis of this equation. The same similarity method can also be applied to some non-linear equations. Separating variables. We can also obtain some solutions of this equation by separating variables. giving us the two ordinary differential equations and solutions of the general form

Intelligence Intensification/Information Sifting. Working in the IT field I come across an unimaginable amount of data on a regular basis in just dealing with day-to-day issues. "Any" task is like that these days - and especially to someone who knows nothing of the task until it's presented. The main thing to remember here is organization. The method you use to organize the data is not important - only that you stick to it for the most part. If something doesn't fit then you either change the item you're having a problem with...or you change your system of organization. Ever watch a flock of birds? One thing you'll notice is that if you try to focus on all the birds at once you'll most likely get confused. If you focus on, say, one bird out of the flock, then you've mentally organized that one bird out of the flock. Guess what...if you keep adding birds out of the flock until you've looked at them all then you've "organized" the flock in your mind. Information in the day-to-day world is no different. Breaking down information into chunks you can deal with is absolutely vital to understanding and remembering. After all, without training you couldn't recall a whole book from memory, right? But you could recall important parts of it. Eventually, if you keep adding to those parts then you would at some point recall the entire book. In another context, this techique has been known for roughly two millennia. Remember the Latin phrase "divide et impera" (divide and rule)? Now, how this relates to sifting is most of the time pretty simple. Organize the data, but focus on the stuff you want. For instance if you want to observe only the red birds in the flock then focus on the color red. Sounds pretty easy right? Well...yes and no. When dealing with large amounts of data you have to do a little estimation of chance up front. What's that mean? Let's say your family bakes a pie with a ring in it for the holiday. Whoever gets the ring wins a prize. You really want that prize so you focus on the ring. Thus you ask some questions of yourself: "Does the ring have any special properties I can use to identify it?" "Does the ring affect the piece of pie it's in in any way?" "Is there anywhere in the pie I "know" the ring is not located?" Suppose the ring is copper and turns the pie green. Then you know that any piece of pie that is the normal color is less likely to have the ring in it. That narrows your focus by a certain amount. Suppose the ring is also pretty big, so the pie bulges where the ring is. Then you know that any flat piece of pie will not have the ring in it. That narrows your focus again. You keep going with this method until you've found the ring. Some things that help this process: "Learning to speed-read. There are plenty of courses out there for this." "Good language skills for the materials you're searching." "Frequent breaks. Spending too much time sifting often blurs your focus." Focus on one thing at a time. You get more done that way. Remember, sunlight spread out only warms things up, but sunlight focused on one place with a magnifying glass can burn holes in things. Know what you want. Know what you are looking for. If you don't know, how will you know if you find it? Delegate. Get people to help you. Parallel searching is helpful, as long as you make sure they know what they are looking for. It's horribly cliche, but big mountains are just piles of pebbles stacked on top of each other. Information is the same... Next. ../Speed Reading/

Field Guide/Chuparosa. The Chuparosa ("hummingbird" in Spanish) is a nectar-rich flowers which attracts hummingbirds, linnets and sparrows. It is said to have been eaten by the Papago tribe.

Guide to Unix/Commands. The Unix command line is often considered difficult to learn. This book aims to help beginners by introducing various commands in lucid and simple language. Unlike most command references, this book is designed to be a self-study guide.

Guide to Unix. This is the Guide to Unix Computing, or for short, Guide to Unix. It describes Unix and Unix-like systems for "users" and "system administrators". It includes guide to Unix commands. Other Wikibooks. __NOEDITSECTION__

Calculus/Introduction to multivariable calculus. This chapter serves as an introduction to multivariable calculus. Multivariable calculus is more complicated than when we were dealing with single-variable functions because more variables means more situations to be concerned about. In the following chapters, we will be discussing limits, differentiation, and integration of multivariable functions, using single-variable calculus as our basis. Topology in R"n". In your previous study of calculus, we have looked at functions and their behavior. Most of these functions we have examined have been all in the form and only occasional examination of functions of two variables. However, the study of functions of "several" variables is quite rich in itself, and has applications in several fields. We write functions of vectors - many variables - as follows: and formula_3 for the function that maps a vector in formula_4 to a vector in formula_5 . Before we can do calculus in formula_5 , we must familiarize ourselves with the structure of formula_5 . We need to know which properties of formula_8 can be extended to formula_5 . This page assumes at least some familiarity with basic linear algebra. Lengths and distances. If we have a vector in formula_10 we can calculate its length using the Pythagorean theorem. For instance, the length of the vector formula_11 is We can generalize this to formula_5 . We define a vector's length, written formula_14 , as the square root of the sum of the squares of each of its components. That is, if we have a vector formula_15 , Now that we have established some concept of length, we can establish the distance between two vectors. We define this distance to be the length of the two vectors' difference. We write this distance formula_17 , and it is This distance function is sometimes referred to as a "metric". Other metrics arise in different circumstances. The metric we have just defined is known as the "Euclidean" metric. Open and closed balls. In formula_8 , we have the concept of an "interval", in that we choose a certain number of other points about some central point. For example, the interval formula_20 is centered about the point 0, and includes points to the left and right of 0. In formula_10 and up, the idea is a little more difficult to carry on. For formula_10 , we need to consider points to the left, right, above, and below a certain point. This may be fine, but for formula_23 we need to include points in more directions. We generalize the idea of the interval by considering all the points that are a given, fixed distance from a certain point - now we know how to calculate distances in formula_5 , we can make our generalization as follows, by introducing the concept of an "open ball" and a "closed ball" respectively, which are analogous to the open and closed interval respectively. In formula_8 , we have seen that the open ball is simply an open interval centered about the point formula_30 . In formula_10 this is a circle with no boundary, and in formula_23 it is a sphere with no outer surface. ("What would the closed ball be?") Boundary points. If we have some area, say a field, then the common sense notion of the "boundary" is the points 'next to' both the inside and outside of the field. For a set, formula_33 , we can define this rigorously by saying the boundary of the set contains all those points such that we can find points both inside and outside the set. We call the set of such points formula_34 . Typically, when it exists the dimension of formula_34 is one lower than the dimension of formula_33 . e.g. the boundary of a volume is a surface and the boundary of a surface is a curve. This isn't always true; but it is true of all the sets we will be using. A set formula_33 is "bounded" if there is some positive number such that we can encompass this set by a closed ball about formula_38 . --> if every point in it is within a finite distance of the origin, i.e there exists some formula_39 such that formula_40 is in S implies formula_41 . Limits. We will focus on the limits of two-variable functions while reviewing the limits of single-variable functions. Multivariable limits are significantly harder than single-variable limits because of different directions. Assume that there is a single-variable function:formula_42In order to ensure that formula_43 exists, we need to test it from two directions: one approaching formula_44 from the left side (formula_45) and the other approaching formula_44 from the right side (formula_47). Recall thatformula_43 exists when formula_49.For example, formula_50 does not exist because formula_51 and formula_52. Now, assume that there is a function with two variables:formula_53If we want to take a limit, for example, formula_54, not only do we need to consider the limit from the direction of the formula_55-axis, we also need to consider the limit from all directions, which includes the formula_56-axis, lines, curves, etc. Generally speaking, if there is one direction where the calculated limit is different from others, the limit does not exist. We will be discussing this in detail here. Differentiable functions. When we expand our scope into the 3-dimensional world, we have significantly more situations to consider. For example, derivatives. In previous chapters, derivatives only have one direction (the formula_55-axis) because there is only one variable.formula_58When we have two or more variables, the rate of change can be calculated in different directions. For example, take a look at the image on the right. This is the graph of a two-variable function. Since there are two variables, the domain will be the whole formula_59-plane. We will graph the output formula_60 on the formula_61-axis. The equation for the function on the right is:formula_62How can we calculate a derivative? The answer is to use partial derivatives. As the name suggests, it can only calculate a derivative "partially" because we can only calculate the rate of change of a graph in one direction. Partial derivatives. Notations are important for partial derivatives.formula_63 means the derivative of formula_60 in the formula_55-axis direction, where we only view the formula_55 as a variable while formula_56 as a constant. formula_68 means the derivative of formula_60 in the formula_56-axis direction, where we only view the formula_56 as a variable while formula_55 as a constant.For simplicity, we will often use various standard abbreviations, so we can write most of the formulae on one line. This can make it easier to see the important details. We can abbreviate partial differentials with a subscript, e.g., formula_73Note that formula_74 in general. They are only sometimes equal. For details, see The_chain_rule_and_Clairaut's_theorem. When we are using a subscript this way we will generally use the Heaviside "D" (which stands for "directional") rather than ∂,formula_75 formula_76 means the derivative of formula_77 in the direction formula_78If we are using subscripts to label the axes, "x1, x2" …, then, rather than having two layers of subscripts, we will use the number as the subscript.formula_79We can also use subscripts for the components of a vector function, formula_80 If we are using subscripts for both the components of a vector and for partial derivatives we will separate them with a comma.formula_81The most widely used notation is formula_82. We will use whichever notation best suits the equation we are working with. Directional derivatives. Normally, a partial derivative of a function with respect to one of its variables, say, "x""j", takes the derivative of that "slice" of that function parallel to the "x""j"'th axis. More precisely, we can think of cutting a function f("x"1...,"x""n") in space along the "x""j"'th axis, with keeping everything but the "x""j" variable constant. From the definition, we have the partial derivative at a point p of the function along this slice as provided this limit exists. Instead of the basis vector, which corresponds to taking the derivative along that axis, we can pick a vector in any direction (which we usually take as being a unit vector), and we take the "directional derivative" of a function as where d is the direction vector. If we want to calculate directional derivatives, calculating them from the limit definition is rather painful, but, we have the following: if f : R"n" → R is differentiable at a point p, |p|=1, There is a closely related formulation which we'll look at in the next section. Gradient vectors. The partial derivatives of a scalar tell us how much it changes if we move along one of the axes. What if we move in a different direction? We'll call the scalar "f", and consider what happens if we move an infinitesimal direction dr=("dx,dy,dz"), using the chain rule. This is the dot product of dr with a vector whose components are the partial derivatives of f, called the gradient of f formula_87 We can form directional derivatives at a point p, in the direction d then by taking the dot product of the gradient with d Notice that grad "f" looks like a vector multiplied by a scalar. This particular combination of partial derivatives is commonplace, so we abbreviate it to We can write the action of taking the gradient vector by writing this as an "operator". Recall that in the one-variable case we can write "d"/"dx" for the action of taking the derivative with respect to "x". This case is similar, but ∇ acts like a vector. We can also write the action of taking the gradient vector as: Properties of the gradient vector. Geometry. For example, if we consider h("x", "y")="x"2+"y"2. The level sets of "h" are concentric circles, centred on the origin, and grad "h" points directly away from the origin, at right angles to the contours. If dr points along the contours of "f", where the function is constant, then "df" will be zero. Since "df" is a dot product, that means that the two vectors, df and grad "f", must be at right angles, i.e. the gradient is at right angles to the contours. Algebraic properties. Like "d/dx", ∇ is linear. For any pair of constants, "a" and "b", and any pair of scalar functions, "f" and "g" Since it's a vector, we can try taking its dot and cross product with other vectors, and with itself. Product and chain rules. Just as with ordinary differentiation, there are product rules for grad, div and curl. We can also write chain rules. In the general case, when both functions are vectors and the composition is defined, we can use the Jacobian defined earlier. where Ju is the Jacobian of u at the point v. Normally J is a matrix but if either the range or the domain of u is R1 then it becomes a vector. In these special cases we can compactly write the chain rule using only vector notation. This substitution can be made in any of the equations containing ∇ Integration. We have already considered differentiation of functions of more than one variable, which leads us to consider how we can meaningfully look at integration. In the single variable case, we interpret the definite integral of a function to mean the area under the function. There is a similar interpretation in the multiple variable case: for example, if we have a paraboloid in R3, we may want to look at the integral of that paraboloid over some region of the "xy" plane, which will be the "volume" under that curve and inside that region. Riemann sums. When looking at these forms of integrals, we look at the Riemann sum. Recall in the one-variable case we divide the interval we are integrating over into rectangles and summing the areas of these rectangles as their widths get smaller and smaller. For the multiple-variable case, we need to do something similar, but the problem arises how to split up R2, or R3, for instance. To do this, we extend the concept of the interval, and consider what we call a "n"-interval. An "n"-interval is a set of points in some rectangular region with sides of some fixed width in each dimension, that is, a set in the form {x∈R"n"|"a"i ≤ "x"i ≤ "b"i with "i" = 0...,"n"}, and its area/size/volume (which we simply call its "measure" to avoid confusion) is the product of the lengths of all its sides. So, an "n"-interval in R2 could be some rectangular partition of the plane, such as {("x","y") | "x" ∈ [0,1] and "y" ∈ [0, 2]|}. Its measure is 2. If we are to consider the Riemann sum now in terms of sub-"n"-intervals of a region Ω, it is where "m"("S""i") is the measure of the division of Ω into "k" sub-"n"-intervals "S""i", and "x"*"i" is a point in "S""i". The index is important - we only perform the sum where "S""i" falls completely within Ω - any "S"i that is not completely contained in Ω we ignore. As we take the limit as "k" goes to infinity, that is, we divide up Ω into finer and finer sub-"n"-intervals, and this sum is the same no matter how we divide up Ω, we get the "integral" of "f" over Ω which we write For two dimensions, we may write and likewise for "n" dimensions. Iterated integrals. Thankfully, we need not always work with Riemann sums every time we want to calculate an integral in more than one variable. There are some results that make life a bit easier for us. For R2, if we have some region bounded between two functions of the other variable (so two functions in the form "f"("x") = "y", or "f"("y") = "x"), between a constant boundary (so, between "x" = "a" and "x" ="b" or "y" = "a" and "y" = "b"), we have An important theorem (called "Fubini's theorem") assures us that this integral is the same as if f is continuous on the domain of integration.

Intelligence Intensification/Proofs. Different levels of proof. There are many different levels of proofs. Some people tend to regard everything as a proof, others nothing. Having a good ability to decide what constitutes a good proof is a clear sign of high intelligence. A person who doesn't know anything about scientific proofs is not able to see through pseudo-scientific proofs and might therefore walk around believing things which are not true. One measure of intelligence is the correspondence of a person's models of reality with reality itself. Elements of a scientific method. The essential elements of a scientific method are iterations and recursions of the following four steps: The above is a hypothetico-deductive method, and includes observation in the first and fourth steps. Each step is subject to peer review for possible mistakes. Characterization. A scientific method depends upon a careful characterization of the subject of the investigation. (The "subject" can also be called "the problem" or the "unknown".) For example, Benjamin Franklin correctly characterized St. Elmo's fire as electrical in nature, but it has taken a long series of experiments and theory to establish this. While seeking the pertinent properties of the subject, this careful thought may also entail some definitions and observations; the observation often demands careful "measurement" and/or counting. The systematic, careful collection of measurements or counts of relevant quantities is often the critical difference between pseudo-sciences, such as alchemy, and a science, such as chemistry. Scientific measurements taken are usually tabulated, graphed, or mapped, and statistical manipulations, such as correlation and regression, performed on them. The measurements might be made in a controlled setting, such as a laboratory, or made on more or less inaccessible or unmanipulatable objects such as stars or human populations. The measurements often require specialized scientific instruments such as thermometers, spectroscopes, or voltmeters, and the progress of a scientific field is usually intimately tied to their invention and development. Measurements demand the use of "operational definitions" of relevant quantities. That is, a scientific quantity is described or defined by how it is measured, as opposed to some more vague, inexact or "idealized" definition. For example, electrical current, measured in amperes, may be operationally defined in terms of the mass of silver deposited in a certain time on an electrode in an electrochemical device that is described in some detail. The operational definition of a thing often relies on comparisons with standards: the operational definition of "mass" ultimately relies on the use of an artifact, such as a certain kilogram of platinum kept in a laboratory in France. The scientific definition of a term sometimes differs substantially from their natural language usage. For example, mass and weight are often used interchangeably in common discourse, but have distinct meanings in physics. Scientific quantities are often characterized by their units of measure which can later be described in terms of conventional physical units when communicating the work. Measurements in scientific work are also usually accompanied by estimates of their uncertainty. The uncertainty is often estimated by making repeated measurements of the desired quantity. Uncertainties may also be calculated by consideration of the uncertainties of the individual underlying quantities that are used. Counts of things, such as the number of people in a nation at a particular time, may also have an uncertainty due to limitations of the method used. Counts may only represent a sample of desired quantities, with an uncertainty that depends upon the sampling method used and the number of samples taken. Hypothesis development. A hypothesis includes a suggested explanation of the subject. It will generally provide a causal explanation or propose some correlation. Observations have the general form of existential statements, stating that some particular instance of the phenomenon being studied has some characteristic. Causal explanations have the general form of universal statements, stating that every instance of the phenomenon has a particular characteristic. It is not deductively valid to infer a universal statement from any series of particular observations. This is the problem of induction. Many solutions to this problem have been suggested, including falsifiability and Bayesian inference. Scientists use whatever they can — their own creativity, ideas from other fields, induction, systematic guessing, etc. — to imagine possible explanations for a phenomenon under study. There are no definitive guidelines for the production of new hypotheses. The history of science is filled with stories of scientists claiming a "flash of inspiration", or a hunch, which then motivated them to look for evidence to support or refute their idea. Michael Polanyi made such creativity the centrepiece of his discussion of methodology. Prediction from the hypothesis. A useful hypothesis will enable predictions, by deductive reasoning, that can be experimentally assessed. If results contradict the predictions, then the hypothesis under test is incorrect or incomplete and requires either revision or abandonment. If results confirm the predictions, then the hypothesis might be correct but is still subject to further testing. Einstein's theory of General Relativity makes several specific predictions about the observable structure of space-time, such as a prediction that light bends in a gravitational field and that the amount of bending depends in a precise way on the strength of that gravitational field. Arthur Eddington's observations made during a 1919 solar eclipse supported General Relativity rather than Newtonian gravitation. "Predictions" refer to experiment designs with a "currently unknown outcome"; the classic example was Edmund Halley's prediction of the year of return of Halley's comet which returned after his death. A prediction (of an unknown) differs from a "consequence" (which can already be known). Once a prediction is made, an experiment is designed to test it. The experiment may seek either confirmation or falsification of the hypothesis. Yet an experiment is not an absolute requirement. In observation based fields of science actual experiments must be designed differently than for the classical laboratory based sciences; for example, the observations of the Chaldeans were utilized in the work of Al-Batani, when he determined a value for the precession of the Earth, in work that spanned thousands of years. Scientists assume an attitude of openness and accountability on the part of those conducting an experiment. Detailed recordkeeping is essential, to aid in recording and reporting on the experimental results, and providing evidence of the effectiveness and integrity of the procedure. They will also assist in reproducing the experimental results. This tradition can be seen in the work of Hipparchus (190 BCE - 120 BCE), when determining a value for the precession of the Earth over 2100 years ago, and 1000 years before Al-Batani. The experiment's integrity should be ascertained by the introduction of a "control". Two virtually identical experiments are run, in only one of which the factor being tested is varied. This serves to further isolate any causal phenomena. For example in testing a drug it is important to carefully test that the supposed effect of the drug is produced only by the drug. Doctors may do this with a double-blind study: two virtually identical groups of w:patients are compared, one of which receives the drug and one of which receives a placebo. Neither the patients nor the doctor know who is getting the real drug, isolating its effects. Once an experiment is complete, a researcher determines whether the results (or data) gathered are what was predicted. If the experimental conclusions fail to match the predictions/hypothesis, then one returns to the failed hypothesis and re-iterates the process. If the experiment(s) appears "successful" - i.e. fits the hypothesis - then its details become published so that others (in theory) may "reproduce" the same experimental results. Testing and improvement. The scientific process is iterative. At any stage it is possible that some consideration will lead the scientist to repeat an earlier part of the process. Failure to develop an interesting hypothesis may lead a scientist to re-define the subject they are considering. Failure of a hypothesis to produce interesting and testable predictions may lead to reconsideration of the hypothesis or of the definition of the subject. Failure of the experiment to produce interesting results may lead the scientist to reconsidering the experimental method, the hypothesis or the definition of the subject. Verification. Science is a social enterprise, and scientific work will become accepted by the community only if they can be verified. Crucially, experimental and theoretical results must be reproduced by others within the science community. Researchers have given their lives for this vision; Georg Wilhelm Richmann was killed by ball lightning to his forehead (1753) when attempting to replicate the 1752 kite experiment of Benjamin Franklin. Reevaluation. "All" scientific knowledge is in a state of flux, for at any time new evidence could be presented that contradicts a long-held hypothesis. A particularly luminous example is the theory of light. Light had long been supposed to be made of particles. Isaac Newton, and before him many of the Classical Greeks, was convinced it was so, but his light-is-particles account was overturned by evidence in favor of a wave theory of light suggested most notably in the early 1800s by Thomas Young, an English physician. Light as waves neatly explained the observed diffraction and interference of light when, to the contrary, the light-as-a-particle theory did not. The wave interpretation of light was widely held to be unassailably correct for most of the 19th century. Around the turn of the century, however, observations were made that a wave theory of light could not explain. This new set of observations could be accounted for by Max Planck's quantum theory (including the photoelectric effect and Brownian motion—both from w:Albert Einstein), but not by a wave theory of light. Nor, for that matter, by the particle theory. More ... Peer review evaluation. Scientific journals use a process of "peer review", in which scientists' manuscripts are submitted by editors of scientific journals to (usually one to three) fellow (usually anonymous) scientists familiar with the field for evaluation. The referees may or may not recommend publication, publication with suggested modifications, or, sometimes, publication in another journal. This serves to keep the scientific literature free of unscientific or crackpot work, helps to cut down on obvious errors, and generally otherwise improve the quality of the scientific literature. Work announced in the popular press before going through this process is generally frowned upon. Sometimes peer review inhibits the circulation of unorthodox work, and at other times may be too permissive. The peer review process is not always successful, but has been very widely adopted by the scientific community. Reproducibility. The reproducibility or replication of scientific observations, while usually described as being very important in a scientific method, is actually seldom actually reported, and is in reality often not done. Referees and editors rightfully and generally reject papers purporting only to reproduce some observations as being unoriginal and not containing anything new. Occasionally reports of a failure to reproduce results are published--mostly in cases where controversy exists or a suspicion of fraud develops. The threat of failure to replicate by others, however, serves as a very effective deterrent for most scientists, who will usually replicate their own data several times before attempting to publish. Evidence and assumptions. Evidence comes in different forms and quality, mostly due to underlying assumptions. An underlying assumption that 'objects heavier than air fall to the ground when dropped' is not likely to incite much disagreement. An underlying assumption like 'aliens abduct humans' however is an extraordinary claim which requires solid proof. Many extraordinary claims also do not survive Occam's razor. Elegance of hypothesis. In evaluating a hypothesis, scientists tend to look for theories that are "elegant" or "beautiful". In contrast to the usual English use of these terms, scientists have more specific meanings in mind. "Elegance" (or "beauty") refers to the ability of a theory to neatly explain as many of the known facts as possible, as simply as possible, or at least in a manner consistent with Occam's Razor while at the same time being aesthetically pleasing. Everyone has reason to learn what constitutes a scientific proof. Even if you never do scientific work, it will help you to evaluate other's work, and to protect yourself against quackery. Maybe even more importantly, it will enable you to think more clearly in general. Statistics. Whenever you hear an advertisement saying a new soap or lotion is scientifically proven to have a positive effect in some sense, statistics have been used (or they lied about the scientificness of the proof). The philosophical ideas behind statistical proofs are these: Suppose, for instance, that you want to see if there is any connection between drinking alcohol during pregnancy and the intelligence of the child. Then you might start with the following: Hypothesis: A mother drinking alcohol during pregnancy "does" lower the intelligence of her child. This gives rise to the following anti-hypothesis or null hypothesis: Null hypothesis: A mother drinking alcohol during pregnancy does not lower the intelligence of her child. Now we want to be 99% sure of our result. That means the risk of error is 1%. After doing a lot of measurements and putting the measurements through the machinery of statistics, we will be able to conclude either: If 2 is the case, we have 'proven' statistically that drinking alcohol during pregnancy lowers the intelligence of the child. Of course this example is stylized. What do we mean by drinking alcohol? What amount, and how regularly? How do we measure intelligence? Those must also be specified. The Axiomatic Method. The axiomatic method is fundamental in every mathematical theory. A complete theory is built of axioms and implications.* Other names for "axiom" are "premise", "postulate", and "assumption". An axiom is always assumed to be true, without discussion, for the sake of argument. Each time we say 'suppose', we describe an axiom. When a statement undoubtedly (logically) follows from another statement, we have an implication. Suppose all human beings have blue eyes. Let that be an axiom. Now suppose that Melinda is a human being. (That's another axiom.) Then Melinda has blue eyes. Any other conclusion about Melinda's eye color would be wrong, because the axioms are defined as true (unless one could prove that the two axioms contradict one another, in which case one would have to be discarded, but that's another story). People have different axioms. Have a look at this belief: "If we don't throw a pancake in the dragons cave each morning, the sun will not rise". You might say that this isn't very logical. But it indeed is, with the right axioms. Axioms: Now, if nobody gives pancakes to the dragon, then the dragon will die (suppose also that the dragon cannot make his own pancakes). But if the dragon is dead, then nobody makes the sun rise! This conclusion is logically derived from the axioms. The point is that even hard core mysticists may use logic in their thinking; only their assumptions are strange. Most people don't consciously think about what are axioms and what are implications when they argue with each other. Also, most people would rather die than change any of their axioms of life. People have the strangest axioms like "Different configurations of the stars have different easily detectable effects on human beings". Maybe they have good reasons to have these axioms. However, having these axioms, they think they are thinking rationally, and they are! As long as the implications follow logically from their axioms, they are thinking rationally! At least, according to one definition of "rational". Another definition might be "a theory is rational if it has a good correlation with physical reality". But then many mathematical theories are not rational; for instance, most non-Euclidian geometries (and ordinary Euclidian geometry too, according to Minkowski-Einstein theory!). And we want mathematical theories to be "rational", so the latter was not a good definition. An intelligent being should be aware of the fact that different people and different cultures have different axioms. It might be a good idea to practice believing in strange things. Be aware of your axioms! Don't believe in them, just regard them as axioms! Change axioms each time you change underwear, if you change your underwear reasonably often. If you have never heard of Occam's Razor, this is the perfect time to learn what it is. It's a principle which roughly says: If you have to choose between two equally good theories for explaining a phenomena, choose the one with the smallest number of axioms. A short repetition: Inductive Reasoning. Induction or inductive reasoning, sometimes called inductive logic, is the process of reasoning in which the conclusion of an argument is very likely to be true, but not certain, given the premises. It is to ascribe properties or relations to types based on limited observations of particular w:tokens; or to formulate laws based on limited observations of recurring phenomenal patterns. Induction is used, for example, in using specific propositions such as: to infer general propositions such as: Validity. Formal logic as most people learn it is deductive rather than inductive. Some philosophers claim to have created systems of inductive logic, but it is controversial whether a logic of induction is even possible. In contrast to deductive reasoning, conclusions arrived at by inductive reasoning do not necessarily have the same degree of certainty as the initial assumptions. For example, a conclusion that all swans are white is obviously wrong, but may have been thought correct in Europe until the settlement of Australia. Inductive arguments are never binding but they may be cogent. Inductive reasoning is deductively invalid. (An argument in formal logic is valid if and only if it is not possible for the premises of the argument to be true whilst the conclusion is false.) The classic philosophical treatment of the problem of induction, meaning the search for a justification for inductive reasoning, was by the Scotsman David Hume. Hume highlighted the fact that our everyday reasoning depends on patterns of repeated experience rather than deductively valid arguments. For example we believe that bread will nourish us because it has in the past, but it is at least conceivable that bread in the future will poison us. Someone who insisted on sound deductive justifications for everything would starve to death, said Hume. Instead of unproductive radical skepticism about everything, he advocated a practical skepticism based on common-sense, where the inevitability of induction is accepted. 20th Century developments have framed the problem of induction very differently. Rather than a choice about what predictions to make about the future, it can be seen as a choice of what concepts to fit to observation (see the entry for grue) or of what graphs to fit to a set of observed data points. Induction is sometimes framed as reasoning about the future from the past, but in its broadest sense it involves reaching conclusions about unobserved things on the basis of what is observed. Inferences about the past from present evidence (e.g. archaeology) count as induction. Induction could also be across space rather than time, e.g. conclusions about the whole universe from what we observe in our galaxy or national economic policy based on past economic performance.

Field Guide/Birds/British Isles. Passeriformes. Thumbnails Accipitriformes. Some classifications also include the Falconidae. Falconiformes. Sometimes included in the Accipitriformes

Australian Government/The Commonwealth Legislature. The legislature of the Commonwealth is the Commonwealth Parliament or Parliament of Australia. The Parliament of Australia is a bicameral parliament consisting of the House of Representatives (the "lower house") and the Senate (the "upper house" or "house of review"). Section 1 of the Constitution of Australia provides that: "The legislative power of the Commonwealth shall be vested in a Federal Parliament, which shall consist of the Queen, a Senate, and a House of Representatives, and which is herein-after called 'The Parliament,' or 'The Parliament of the Commonwealth'." Queen Elizabeth II is, in her capacity as Queen of Australia, Australia's head of state, but her constitutional functions in Australia are delegated to the Governor-General of Australia. The House of Representatives consists of 150 members elected from single-member constituencies of approximately equal population. The Senate consists of 76 members: 12 Senators are elected from each of the six states and two from each of the two territories. The principal function of the Parliament is to pass laws, or legislation. Any Member or Senator may introduce a proposed law (a bill), except for a money bill (a bill proposing an expenditure or levying a tax), which must be introduced in the House of Representatives. In practice, the great majority of bills are introduced by ministers. Bills introduced by other Members are called private members' bills. All bills must be passed by both Houses to become law. The Senate has the same legislative powers as the House, except that it may not amend money bills, only pass or reject them. The Parliament performs other functions besides legislation. It can discuss urgency motions or matters of public importance: these provide a forum for debates on public policy matters. Members can move motions of censure against the government or against individual ministers. On most sitting days in both Houses there is a session called Question Time at which Members and Senators address questions to the Prime Minister and other ministers. Members and Senators can also present petitions from their constituents. Both Houses have an extensive system of committees in which draft bills are debated, evidence is taken and public servants are questioned. Members of the Australian Parliament do not have legal immunity: they can be arrested and tried for any offence. They do, however, have Parliamentary privilege: they cannot be sued for anything they say about each other or about persons outside the Parliament. This privilege extends to reporting in the media of anything a Member or Senator says. There is a legal offence called contempt of Parliament. A person who speaks or acts in a manner contemptuous of the Parliament or its members can be tried and, if convicted, imprisoned. The Parliament used to have the power of hearing such cases itself, and did so in the Browne-Fitzpatrick case of 1955. This power has now been delegated to the courts, but no-one has been prosecuted for this offence. The House of Representatives. The 150 members of the house are elected from single-member geographic districts (popularly known as "seats" but officially known as "Commonwealth Electoral Divisions") which are intended to represent reasonably contiguous regions, with relatively equal population in each of about 80 000 people. Voting is by the preferential system. According to Australia's Constitution, the powers of both houses are nearly equal with the consent of both houses needed to pass legislation. In practice, however, the Lower House is far stronger in some ways, and far weaker in others. By convention, the party or coalition in the lower house with a majority is invited by the Governor-General to form government, and thus the leader of the party in the lower house becomes the Prime Minister of Australia and his senior colleagues ministers responsible for various government departments. Bills appropriating money can also only be introduced or modified in the lower house. Thus, only parties in the lower house can govern. However, in the rigid Australian party system, this ensures that virtually all contentious votes are along party lines, and the government always has a majority in those votes. The Opposition's only real role in the House is to present arguments why the government's policies and legislation are wrong, and attempt to embarrass the government as much as possible by asking difficult questions at question time. In a reflection of the color scheme of the United Kingdom House of Commons, the House of Representatives is decorated in green. The Senate. The voting system for the Senate has changed twice since it was created. The original arrangement was a first past the post block voting mechanism. In 1919, it was changed to preferential block voting. Block voting tended to grant landslide majorities very easily. In 1946, the Australian Labor Party government won 30 out of the 33 Senate seats. In 1948, partially in response to this extreme situation, they introduced proportional representation in the Senate. From a comparative governmental perspective, the Australian Senate is almost unique in that unlike the upper house in other Westminster system governments, the Senate is not a vestigial body with limited legislative power but rather plays and is intended to play an active role in legislation. Rather than being modelled after the House of Lords the Australian Senate was in part modelled after the United States Senate and was intended to give small rural states added voice in a Federal legislature, while also fulfilling the revising role of an upper house in the Westminster system. Although the Prime Minister is answerable to, and selected from the House of Representatives (the "lower house"), other ministers are drawn from either house and the two houses have almost equal legislative power. As with most upper chambers in bicameral parliaments, it cannot introduce Appropriation Bills or impose taxation, that role being reserved for the universally elected lower chamber. That degree of equality between the Australian Senate and House of Representatives is in part due to the age of the Australian constitution; it was enacted before the confrontation in 1909 between the British House of Commons and House of Lords, over estate taxes, which ultimately resulted in the restrictions in the powers of the House of Lords in the Parliament Act. The Senate thus reflected the pre-1911 Lords-Commons relationship, namely a house with theoretically wide powers that are by convention not widely used, but which in a crisis could be. However in one area it possesses a highly sensitive power, namely the right to withdraw or block Supply, i.e. government control of exchequer funding. In most democracies, this power does not exist within upper houses. In the Australian case, probably due to its age and its pre-dating of the British Parliament Act, 1911, that power, once possessed by the House of Lords, is still possessed by the Australian Senate. What it means is that in strict constitutional terms, the Government is also answerable to the Australian Senate, given that the loss of Supply in parliamentary democracies automatically requires the resignation of a government or ministry, or alternatively the calling of a general election, because without access to exchequer funding a government cannot function and would face bankruptcy. However, as in the United Kingdom prior to the 1909 clash over David Lloyd George's budget between the House of Commons and House of Lords, a general convention built up in which the upper house would not use its power to block exchequer funding, leaving that responsibility to the democratically representative lower house. Indeed to break convention in this area is sometimes described as a parliamentary nuclear option because of its political, financial and governmental impact. In 1975, in a dispute strikingly similar to Britain's 1909 clash (an upper House breaking convention by withdrawing Supply on the basis that it was reacting to a breach of another fundamental convention by the government - which the government denied) the Senate did refuse to pass a required financial measure, producing a 'nuclear'-style crisis which resulted in a stand off, a decision of the Prime Minister not to resign or seek a dissolution, and the eventual intervention of the Governor-General to withdraw the commission of the Australian Prime Minister (in effect dismissing him) and the appointment of a minority government from the opposition in the House of Representatives and the calling of a general election. In practice, however, most legislation (except for "Private Member's Bills") in the Australian Parliament is initiated by the Government, which has control over the lower house. It is then passed to the Senate, which may amend the bill or refuse to pass it. In the majority of cases, voting is along party lines.

Australian Government/The Commonwealth Constitution. < Australian Government/Contents The Australian Constitution is actually the Commonwealth of Australia Constitution Act 1900, which was passed through the United Kingdom Parliament. It provided the new system of government for the new federation, which consisted at its inception on 1 January 1901 of the former separate colonies of New South Wales, Victoria, South Australia, Queensland, Tasmania and Western Australia. Head of State. The Act vested authority in the Queen, making her the Australian head of state similar to other Commonwealth Realms. In 1973, the monarch was formally designated as 'Queen of Australia'. A representative of the Queen was provided for, known as the Governor-General, who in practice fulfils most of the roles normally possessed by a head of state. Some consider the Governor-General of Australia to be the de facto head of state as the British monarch rarely exercises the reserve powers that the constitution grants to the Crown; however the constitution makes clear that the Governor-General is in no sense a head of state, merely a head of state's representative who in the name of the head of state, or in his own name as representative of the head of state, carries out specified functions and exercises certain powers. Parliament. Section 1 (of Chapter I) provided that the legislative power was to be vested in Federal parliament, known as the 'Parliament of the Commonwealth', consisting of the Queen, an upper house, called the Senate, and a lower house, called the 'House of Representatives. Executive Authority. According to Section 61 (of Chapter II): "the executive power of the Commonweath is vested in the Queen and is exercisable by the Governor-General as the Queen's representative, and extends to the execution and maintenance of this Constitution, and the laws of the Commonweath." Article 62 provided for a Federal Executive Council to 'advise' the Governor-General in the governance of the Commonwealth. Though the language indicated that the Executive Council was answerable to the Governor-General, in reality it is answerable to the House of Representatives, though the fact that the Senate possesses the power to withdraw Supply complicates the situation, given that loss of Supply in parliamentary democracies has the most severe implications for a government, given that it in theory should either resign or seek a parliamentary dissolution, should Supply be lost or not granted. "See:The Commonwealth Executive". The Judiciary. The judicial power of the Commonwealth was vested by Section 71 of Chapter III in a federal supreme court to be called the High Court of Australia. It was to be presided over by a Chief Justice. "See:The Judiciary" The States. Section 106 of Chapter V provided for the continuation of the constitutions of the various states, subject to the provisions of the federal constitution. Amendments to the Constitution. Section 128 of Chapter VIII provided that constitutional amendments required and

Guide to Unix/Files. /etc/. /etc/fstab. The fstab (for "file systems table") file is commonly found on Unix and Unix-like systems and is part of the system configuration. The fstab file typically lists all used disks and disk partitions, and indicates how they are to be used or otherwise integrated into the overall system's file system. Traditionally, the fstab was only read by programs, and not written to. However, more modern system administration tools can automatically build and edit fstab, or act as graphical editors for it. It is the duty of the system administrator to properly create and maintain this file. The file may have other names on a given Unix variant; for example, it is codice_1 on Solaris. Example The following is an example of a fstab file on a Red Hat Linux system: The first column indicates the device name or other means of locating the partition or data source. The second column indicates where the data is to be attached to the filesystem. The third column indicates the filesystem type, or algorithm to use to interpret the filesystem. The fourth column gives options, including if the filesystem should be mounted at boot. The fifth column adjusts the archiving schedule for the partition (used by dump). The sixth column indicates the order in which the fsck utility will scan the partitions for errors when the computer powers on. A value of zero in either of the last 2 columns disables the corresponding feature (http://www.humbug.org.au/talks/fstab/fstab_structure.html). To get more information about the fstab file you can read the man page about it. The Kfstab graphical configuration utility is available for KDE for editing fstab. See also /etc/group. /etc/group stores the definitive list of the users groups and their members. A typical entry is: root::0:root,alice It has four sections which going from left to right are, /etc/passwd. /etc/passwd is the user authentication database, it contains a list of users and their associated internal user id numbers. Historically it also included passwords, however as this file needs to world readable (so all programs can use it to convert between username and user id) it is no longer considered secure to keep passwords in this file. An entry in this file is of the form: alice:*:134:20:Alice Monkey:/home/alice/:/bin/bash It has seven sections which going from left to right are, /etc/profile. /etc/profile contains the system default settings for users who login using the Bourne shell, "/bin/sh". When these users login, the Bourne shell runs the commands in this file before giving the shell prompt to the user. Most of these commands are variable assignments which configure the behavior of the shell. Some Bourne-compatible shells also use this file, but other shells, such as the C shell, do not. /etc/shadow. /etc/shadow contains the passwords for users in systems which use shadowing. alice:43SrweDe3F:621:5:30:10:100:900: The sections are: /etc/sysctl.conf. /etc/sysctl.conf configures the behavior of the running Unix kernel. During system boot, the scripts read this file and use "sysctl" to set the parameters shown in the file. Changing the file has no effect before the next reboot. /dev/. /dev/cdrom. /dev/cdrom is not an actual device, but on many systems it is a symbolic link to the actual CD device. For example, a Linux system with /dev/hdb for its floppy drive is likely to have a link /dev/cdrom which redirects to /dev/hdb. /dev/fd*. At Linux, /dev/fd0 is the first floppy disk drive at the system. Use /dev/fd0H1440 to operate the first floppy drive in high density mode. Generally, this is invoked when formatting a floppy drive for a particular density. Slackware comes with drivers that allow for formatting a 3.5" diskette with up to 1.7MB of space. Red Hat and Mandrake do not contain these device driver files by default. Likewise, /dev/fd1 is the second floppy disk drive. /dev/hd*. At Linux, /dev/hda is the first IDE hard drive. The second drive is either /dev/hdb or /dev/hdc, depending on the hardware configuration. Some IDE hardware allows up to four drives, including /dev/hdd. Many machines have one hard drive (hda) and one cdrom drive (hdc on many machines, but hdb on some). Often, /dev/cdrom is a symbolic link to the cdrom drive. Partitions are numbered from 1, like /dev/hda1, /dev/hda2, ... /dev/null. /dev/null is a do-nothing device to use when one wants to ignore or delete program output. This file is useful when a program expects to save to a file, but you want not to save anything. This file can also be used as input to a program to represent an empty file. There is no actual hardware associated with the /dev/null device. "Examples:" Deleting file called "x" (command rm x) sometimes causes an error, for example if the file does not exist: $ rm x rm: x: No such file or directory One can hide the error by redirecting it to a file. By using /dev/null as the file, the error never saves to an actual file. "Bourne shell:" $ rm x > /dev/null 2>&1 In the Bourne shell, the "2>&1" redirects the standard error of "rm" (where the error appears) to standard output, then the ">" redirects the standard output to /dev/null. One way to make an empty file called "y" is: $ cat /dev/null > y The "cat" command copies the file "/dev/null" to standard output, and the shell operator ">" redirects this output to "y". The "/dev/null" file seems empty when read, so the file "y" appears, but is also empty. (Note that in this case, simply "> y" will do the same thing.) Dot files. There is some redundancy across these programs. For example, the look and behavior of emacs can be customized by using the .emacs file, but also by adding the appropriate modifications to the .Xdefaults file. Default versions of these files are often installed in users' home directories when the software packages that use them are installed. If a program doesn't find its configuration file in the user's home directory, it will often fall back on a system-wide default configuration file installed in one of the subdirectories that the package lives in. Directories. Different distributions have different directory structures, despite attempts at standardization such as the Linux Filesystem Hierarchy Standard (FHS) organization. External link: Modified Directory Structure

Guitar. The aim of this book is to introduce beginners to the basic concepts of the guitar. Important techniques are given their own sections with exercises and examples provided. It includes an extensive section on amplifiers and numerous tips and advice on all aspects of the guitar. Music theory concepts are clearly presented and explained. Links. __NOEDITSECTION__

Guitar/Anatomy of a Guitar. This chapter presents an overview of the different parts most commonly found on the three main types of guitar. On Acoustics and Electrics. Body. The body of a guitar consists of a treble or upper bout (the half of the guitar closest to the neck), the bass or lower bout (the wider half of the guitar), and the waist bout (the narrow section between the treble and bass bouts). The body is one of the most important factors in shaping the overall tone of a guitar. It provides the resonance that shapes the tonal qualities. It determines the volume of acoustic guitars and affects the sustain of electric guitars. Resonance is affected by: The woods listed below are used in the construction of both acoustic and electric guitars. Body top. Some electric guitars have an extra top added to the body to blend the tonal qualities of different types of wood together. Maple with figuring is a popular top and produces a pronounced look and tone (adds brightness). Body tops are not used on acoustics since the layering of two pieces of wood for the table would inhibit the resonance and dull the tone. Bridge. The bridge is found on the lower bout of the body and forms one one end of the vibrating length of the strings, the other end being the "nut", which is located at the end of the fretboard. On most electric guitars, the height of the bridge is adjustable by screws, allowing the guitars "action" – distance that the strings sit above the frets – to be raised or lowered. On most electrics, the horizontal position of the bridge saddles is also screw-adjustable, allowing the guitar's intonation to be set by the user. These adjustments are discussed further in the Adjusting the Guitar section. Depending on the guitar, the strings may terminate at the bridge or just pass over it. The bridge of an acoustic consists of two parts: a saddle and the tie block. Saddles are either a piece of plastic or polished bone. Acoustic guitar saddles are made with a smooth top edge (no notches) and the base of the saddle is seated in a groove cut into the tie block. The wood tie block of a classical guitar is glued to the lower bout and acts as a string terminator. A classical guitar string is pushed through the hole in the tie block and the string is then brought back under itself three or four times and pulled tight to form a knot. Once the saddle is seated in the groove of the tie block the tension of the strings clamp it. Steel string acoustics also have a saddle and tie block though due to the strings having terminating end balls there is no need to knot. The ball end of the string is thrust into a vertical hole in the tie block and secured with a pin. The position and height of the saddle and tie block on acoustics are set by the manufacturer, and usually do not require user adjustment. Adjustments to the height of the acoustic saddle are possible by shaving (lowering) the saddle though this job is best left to a luthier since if too much is taken off, a new saddle will required unless one resorts to the temporizing measure of shims. The design of bridges varies between manufacturers and the above generic descriptions may not apply to some guitars. Fretboard and Frets. The fretboard is a piece of wood that is glued to the front of the neck. These are commonly made of rosewood though other hard woods such as ebony may also be used. Embedded in the fretboard are a number of metal "frets" (fret-wire) usually numbering twenty one to twenty-four. Strings are pressed down behind a fret which changes the length that is left free to vibrate thereby producing a different note. A simple demonstration is to be found on the twelfth fret. On all guitars this is the fret that divides the string exactly in half and produces a note an octave higher than the open note. Any open string that maintains its original tension and is halved produces its octave. This applies to all stringed instruments including the piano and violin. There are a variety of fret designs. Jumbo frets are higher and wider than normal frets and require less fretboard contact to sound a clear note. Medium frets are closer to the board and must be firmly in contact with the fretboard to sound a clear note. Some guitarist prefer jumbo frets due to the ease with which you can bend strings and the faster play offered by less fretboard contact. As with many design elements of the guitar this is a subjective area that is more personal preference rather than advantage. Good technique is not dependent on fret size The first fret is the one nearest the nut. Some manufacturers place a zero fret immediately after the nut and the strings sit on the zero fret. This brings the sound of the open strings nearer to the quality of a fretted note. A fretboard may have decorative inlays at the 3rd, 5th, 7th, 9th and 12th frets which serve as markers for the positions of the guitar. Fretboard inlays can be highly decorative or simple shapes and on expensive guitars may be made from exotic material like mother of pearl or abalone. Head. The headstock lies at the end of the guitar's neck. The purpose of the headstock is to support the tuners, which terminates the strings of the instrument. The tuners are attached to tuning pegs and this allows the guitarist to lower or raise the pitch of the string. A secondary purpose of the headstock is identification and many guitar manufacturers choose to use a distinctive headstock shape often in combination with the name of the model and a trademark logo. On some guitars the model name and trademark logo may be created using inlaid materials though decals are also commonly used. Neck. The neck can be a single piece of wood or several pieces glued together and cut to shape. The fretboard is a separate piece of wood that is attached to the neck. Necks can be glued to the body (set neck) or bolted on. Set necks are usually found on acoustic guitars and many other instruments including the violin, lute and cello. The bolt-on neck is a design feature more commonly associated with electric guitars. Most necks are wood though alternative materials such as carbon fibre composites have been used. Nuts. All strings pass through a "nut" at the headstock end of the fretboard. Its function is to maintain correct string spacing and alignment so that the strings feed into their respective tuning pegs. On acoustic guitars the nut and saddle are usually made of the same material. Electric guitars commonly use plastic or synthetic nuts though sometimes metal is used. As tremolo bars can cause tuning problems, guitars equipped with them usually have some manner of locking nut, where the strings are clamped down. Tip: Some guitarists lubricate the nut grooves so that the strings move more smoothly. You can do this at home with a soft graphite pencil. There's no need for excessive marking with the pencil just a few swipes through the groove should deposit enough graphite. Pick Guard. Guitars, in common with all wood instruments, are prone to dents, scratches and wear. A pick guard (also known as a scratch plate) protects the body of the guitar at the point of most contact. Some electric guitars have raised pick guards so your pick is directed out and away from the pots and strings. Pick guards sometimes need replacing due to wear or damage. In the case of an expensive or rare guitar, which may have a tortoise shell pick guard, the guitar will have to be sent to an experienced luthier. Sound Hole. Sound holes are found on all acoustics. Their purpose is to allow the air pressure to stay equalized so that the soundboard can vibrate. Archtop guitars have f-shaped sound holes - a design feature they share with the violin, viola, cello and double bass. Round sound holes usually have a decorated edge based on a geometrical design known as a rosette. On modern guitars these decorations are machine-made though some luthiers of expensive guitars still use the traditional method of laying by hand small pieces of exotic material like mother of pearl. Truss Rod. Steel-strung guitars, whether acoustic or electric, have a metal truss rod that runs the length of the neck under the fretboard. Strengthening the neck with a truss rod counteracts the tension exerted by the strings and allows the curvature of the neck to be adjusted. Classical guitars do not require a truss rod due to the lower tension of nylon strings. Some less expensive steel-string acoustics do not have a truss rod. Adjusting the truss rod can have a marked impact on the tuning and playability of a guitar. Truss rod adjustments should be performed with great care as it is possible to damage the neck or the truss rod. Overtightening is especially to be avoided because if the truss rod snaps or its threads strip, a very involved repair job will be necessitated. For this reason, the truss rod nut should only be tightened one quarter turn at a time. It should then be left alone for a day so the neck has time to assume the new shape. If further adjustment is necessary, repeat the process – a quarter turn, max. If one is not mechanically inclined, truss rod adjustment should be deferred to someone who is, including perhaps a guitar repair-person. Tuning Pegs. The strings are tensioned by means of tuning pegs, also known as tuners, tuning machines, or machine heads. A high gear ratio gives smoother, easier tuning but is not essential. A guitar can be tuned successfully even with fairly cheap tuners. Most guitar tuning problems originate in other areas, especially string sticking in the nut slots, and dirty or worn strings. Due to the tension of the strings and the constant turning of the pegs the screws that secure the tuners to the headstock may loosen. It is recommended that you check that they are screwed in tightly though avoid over-tightening which may cause alignment problems or damage the screw head or the wood. Electric Guitar. Pickup. A pickup is a magnet wrapped in a coil of wire. When a string is plucked the vibration of the string causes the magnetic flux to vary, inducing a voltage in the coil. This is then amplified and played through a speaker. "Humbucking pickups" use two coils having opposite polarity, placed in magnetic fields of opposite direction. The result is that the signal from the strings is doubled, while electromagnetic noise from other sources is cancelled. Thus humbucking pickups, or "humbuckers", have less hum and other objectionable noise pickup than single-coil pickups. However, having two coils in series, their inductance is generally higher, which decreases their high-frequency response so they sound less bright than single-coils. Humbuckers are also less able to reproduce very high harmonic frequencies of the string because their magnetic field is less localised than the field of a single coil pickup. "Active pickups" are pickups with a built in preamplifier. This can reduce noise pickup and reduce signal losses in the guitar cord. Against these advantages must be weighed cost and the fact that a power source must be provided for the preamplifier. The power source can be a 9 volt battery on the guitar, or "phantom power" piped in via the guitar cord. Passive single coils are the standard pickup for Fender Stratocasters. They have a bright and twangy clean sound but traditionally have a relatively low output voltge. They are susceptible to noise pickup from transformers – including the transformer in your amp –, neon signs, and other AC electrical equipment. This may or may not be a critical problem, depending on your physical location, and on whether you use a clean sound or a sound with a lot of amplifier compression/distortion which will make electromagnetic noise pickup more noticeable and more of a problem. Some single coils, such as the P-90, are larger than regular single coils, and thus warmer than a standard single coil. However they still retain more of a single coil sound and still can pick-up background hum. Another single coil design is the Lipstick, commonly found on Danelectro Guitars, where the entire pickup is placed in a metal enclosure with a small gap left between the two metal halves. Lipstick pickups tend to be bright sounding and the metal case results in some reduction of hum. Active single-coils can have higher output and enhanced sensitivity. Humbuckers provide a warm and fat sound. They typically have a higher output voltage than single-coils which suits them to overdriving an amplifier so as to create a heavy sound with lots of sustain and distortion. Some humbuckers allow coil tap (using only one of the coils) or parallel connection which provides a sound similar to a single coil. Pickups of every type are found in all genres of music, although single-coils, with their bright, twangy sound are often associated with country; and humbuckers producing thick, warm sounds with the rhythm playing in rock. In heavy metal music, humbuckers are the norm because of the necessity of having a low-noise pickup when the amplifier gain is so high. Pickup Arrangements. There are many different arrangements for pickups. The most basic is a single pickup near the bridge. Pickup Selector. Every electric guitar, except those with a single pickup, has a pickup selector. Guitars with two pickups have a three-way switch which allows the guitarist to select either the neck pickup or the bridge pickup. When the switch is in the middle position both pickups are used. On guitars with three pickups there is usually a five way switch. The positions are: Tremolo Bar. A tremolo bar alters the pitch of the strings. Pushing down on the bar lowers the pitch of the strings and pulling up raises the pitch. Rapidly pushing and releasing will produce a modulation in pitch called "vibrato". Vibrato is often confused with tremolo, a volume modulation effect found on amplifiers, hence the misnomer tremolo bar. Originally used just for vibrato; the modern improvements in guitar design, amplifiers and effects has allowed guitarists to create a new palette of tremolo bar sound effects like the popular "dive bomb". There are four kinds of tremolo: Pots. Almost all electric guitars have at least two pots (potentiometers) to control the volume and tone. The usual tone control circuit uses the inductance of the pickup's coil in conjunction with an additional capacitor to reduce high frequencies when the pot is turned counterclockwise. Near the minimum position (at about 2 or 3 out of 10) there is usually a noticeable midrange peak or nasality resulting from the capacitor and inductor resonating with each other. This is a distinctive sound not easily achieved with the tone controls of the amplifier alone. A Fender Stratocaster will typically have one master volume pot and two tone pots for the neck pickup and middle pickup. Electric Guitar Necks. This section describes the different methods used for attaching the neck to an electric guitar: Bolt-on neck - the neck is attached to the body with bolts which are held by a mounting plate for increased stability. The mounting plate can make accessing the higher frets difficult so some manufacturers, notably Ibanez, use a hidden plate. The bolt-on neck is a standard design used by Fender. Set neck - the neck is attached to the body with adhesive. This is the method used on acoustics and rarely is it used for mass-produced electrics. Electric guitars that feature a set neck have to be built to a high standard since once glued on the neck is permanent and cannot be adjusted. Set necks are commonly found on more costly electric guitars. Gibson and Epiphone use set necks which is claimed to have these advantages over a bolt-on neck: Thru-body neck - the neck extends the entire length of the body. The strings, fretboard, pickups and bridge are all mounted on the thru-body neck. The ears or wings (the bouts) are attached or glued to the central stick. The wings may be book-matched in order to give a symmetrical appearance. The thru-body neck is usually found on high-end guitars since the design is not favoured by mass-production manufacturers. It is more common on basses than guitars. The thru-body neck allows easier access to the higher frets because there is no heel and is considered by some guitarists to offer greater sustain.

Biochemistry/Proteins/Introduction. Protein role and importance. Proteins are among the fundamental molecules of biology. They are common to all life present on Earth today, and are responsible for most of the complex functions that make life possible. They are also the major structural constituent of living beings. According to the /Central Dogma/ of Molecular Biology (proposed by Francis Crick in 1958), information is transferred from DNA to RNA to proteins. DNA functions as a storage medium for the information necessary to synthesize proteins, and RNA is responsible for (among other things) the translation of this information into protein molecules, as part of the ribosome. Virtually all the complex chemical functions of the living cell are performed by protein-based catalysts called enzymes. Specifically, enzymes either make or break chemical bonds. Protein enzymes should not be confused with RNA-based enzymes (also called ribozymes), a group of macromolecules that perform functions similar to protein enzymes. Further, most of the scaffolding that holds cells and organelles together is made of proteins. In addition to their catalytic functions, proteins can transmit and commute signals from the extracellular environment, duplicate genetic information, assist in transforming the energy in light and chemicals with astonishing efficiency, convert chemical energy into mechanical work, and carry molecules between cell compartments. Functions not performed by proteins. Proteins do so much that it's important to note what proteins don't do. Currently there are no known proteins that can directly replicate themselves. Prions are no exception to this rule. It is theorized that prions may be able to act as a structural template for other chemically (but not structurally) identical proteins, but they can't function as a native template for protein synthesis "de novo". Proteins don't act as fundamental energy reserves in most organisms, as their metabolism is slower and inefficient compared to sugars or lipids. They are, on the other hand, a fundamental nitrogen and amino acid reserve for many organisms. Proteins do not directly function as a membrane in most organisms, except viruses; however, they are often important components of these structures, lending both stability and structural support. =Proteins as polymers of amino acids= Composition and Features. Proteins are composed of a linear (not branched and not forming rings) polymer of amino acids. The twenty genetically encoded amino acids are molecules that share a central core: The α-carbon is bonded to a primary amino (-NH2) terminus, a carboxylic acid (-COOH) terminus, a hydrogen atom, and the amino acid side chain, also called the "R-group". The R-group determines the identity of the amino acid. In an aqueous solution, at physiological pH (~6.8), the amino group will be in the protonated -NH3+ form, and the carboxylic acid will be in the deprotonated -COO- form, forming a zwitterion. Most amino acids that make up proteins are L-isomers, although a few exotic creatures use D-isomers in their proteins. It is important to note that levorotatory (L) and dextrorotatory (D) are not specific to rectus (R) and sinister (S) configurations. A levorotatory form of a protein can be either R or S configuration. Levorotatory and dextrorotatory refer to how the proteins bend light in a polarimeter. Amino acids polymerize via peptide bonds, which is a type of amide bond. A peptide bond is formed upon the dehydration of the carboxy-terminus of one amino acid with the amine terminus of a second amino acid. The resulting carbonyl group's carbon atom is directly bound to the nitrogen atom of a secondary amine. A peptide chain will have an unbound amino group free at one end (called the N-terminus) and a single free carboxylate group at the other end (called the C-terminus). The written list of the amino acids linked together in a protein, in order, is called its primary structure. By convention, peptide sequences are written from N-terminus to C-terminus. This convention mimics the way polypeptides are synthesized by the ribosome in the cell. Small polymers of less than 20 amino acids long are more often called peptides or polypeptides. Proteins can have sequences as short as 20-30 amino acids to gigantic molecules of more than 3,000 amino acids (like Titin,a human muscle protein). Genetically-Encoded Amino Acids. While there are theoretically billions of possible amino acids, most proteins are formed of only 20 amino acids: the genetically-encoded (or more precisely, "proteogenic") amino acids. Note that all amino acids except glycine have a chiral center at their α-carbons. (Glycine has two hydrogens on its α-carbon, and therefore it is achiral.) Besides glycine, all proteogenic amino acids are -amino acids, meaning they have the same absolute configuration as -glyceraldehyde. This is the same as the S-configuration, with the exception of cysteine, which contains a sulfur atom in its side chain, and so the naming priority changes. -amino acids are sometimes found in nature, as in the cell walls of certain bacteria, but they are rarely incorporated into protein chains. The side-chains of proteogenic amino acids are quite varied: they range from a single hydrogen atom (as for glycine, the simplest amino acid) to the indole heterocycle, as found in tryptophan. There are polar, charged and hydrophobic amino acids. The chemical richness of amino acids is at the base of the complexity and versatility of proteins. Post-Translational Modification. There are two major types of post-translational modifications to protein: those that cleave the bonds of the peptide backbone and those that add or remove functional groups to the sidechains of individual amino acids. In the first type of post-translational modification, specialized enzymes called proteases recognize specific amino acids of a protein and break the associated peptide bond, thereby irreversibly modifying the primary structure. In the second major type of modification, the amino acid side chains of a given protein are chemically modified by enzymatic reactions or are spontaneously formed (non-enzymatic). Examples of sidechain modifications are quite numerous, but common ones include oxidation, acylation, glycosylation (addition of a glycan, or sugar), methylation, and phosphorylation. Both types of post-translational modifications are capable of exerting positive and negative control over a given protein or enzyme. The importance of protein structure. Generally speaking, the function of a protein is completely determined by its structure. Molecules like DNA, which perform a fairly small set of functions, have an almost fixed structure that's fairly independent from sequence. By contrast, protein molecules perform functions as different as digesting sugars or moving muscles. To perform so many different functions, proteins come in many different structures. The protein function is almost completely dependent on protein structure. Enzymes must recognize and react with their substrates with precise positioning of critical chemical groups in the three-dimensional space. Scaffold proteins must be able to precisely dock other proteins or components and position them in space in the correct fashion. Structural proteins like Collagen must face mechanical stresses and be able to build a regular matrix where cells can adhere and proliferate. Motor proteins must reversibly convert chemical energy in movement, in a precise fashion. Protein folding depends on sequence. As Anfinsen demonstrated in the 1960's, proteins acquire their structure by spontaneous folding of the polypeptide chain into the minimal energy configuration. Most proteins require no external factors in order to fold (although specialized protein exist in cells, called chaperones, that help other, misfolded, proteins acquire their correct structure) — the protein sequence itself uniquely determines the structure. Often the whole process takes place in milliseconds. Despite the apparent chemical simplicity of proteins, the vast number of permutations of twenty amino acids in a linear sequence leads to an amazing number of different protein folds. Nevertheless, protein structures share several characteristics in common: they are almost all built of a few secondary structure elements (short-range structural patterns that are recurrent in protein structures) and even the way these elements combine is often repeated in common motifs. Nonetheless, it is still impossible to know what structure a given protein sequence will yield in solution. This is known as the protein folding problem, and it is one of the most important open problems in modern molecular biology. Protein denaturation. Proteins can lose their structure if put in unsuitable chemical (e.g. high or low pH; high salt concentrations; hydrophobic environment) or physical (e.g. high temperature; high pressure) conditions. This process is called denaturation. Denatured proteins have no defined structure and, especially if concentrated, tend to aggregate into insoluble masses. Protein denaturation is by no means an exotic event: a boiled egg becomes solid just because of denaturation and subsequent aggregation of its proteins. Denatured proteins can sometimes refold when put again in the correct environment, but sometimes the process is irreversible (especially after aggregation: the boiled egg is again an example). It is finally the proteins which are responsible for susceptibility or resistance to a pathogen or parasite. Proteins can fold into domains. A significant number of proteins, especially large proteins, have a structure divided into several independent domains. These domains can often perform specific functions in a protein. For example, a cell membrane receptor might have an extracellular domain to bind a target molecule and an intracellular domain that binds other proteins inside the cell, thereby transducing a signal across the cell membrane. The domain of a protein is determined by the secondary structure of a protein there are four main types of domain structures: alpha-helix, beta-sheet, beta-turns, and random coil. The alpha-helix is when the polypeptide chain forms a helix shape with the amino acids side chains sticking out, usually about 10 amino acids long. The alpha-helix gets its strength by forming internal hydrogen bonds, that occur between amino acid 1 and 4 along the length of the helix. A high concentration of Glycine's in a row tend towards the alpha-helix conformation. The beta-sheet structure is composed of poly-peptide chains stacking forming hydrogen bonds between the sheets. You can form parallel sheets by stacking in the same direction N-C terminal on top of N-C terminal or form anti-parallel by stacking in opposite directions N-C terminal on top of C-N terminal. Beta-turns link two anti-parallel beta strands by a 4 amino acid loop in a defined conformation. A random coil is a portion of the protein that has no defined secondary structure. Domains of a protein then come from unique portions of the peptide that are made up of these types of secondary structure.

Biochemistry/Proteins/The chemistry of proteins. Amino acid structure and chemistry. General features. Amino acids consist of a primary amine bound to an aliphatic carbon atom (the so-called α-carbon), which in turn is bound to a carboxylic acid group. At least one hydrogen atom is bound to the α-carbon; in addition, the α-carbon bears a side chain, which is different for different amino acids. In a neutral aqueous solution, amino acids exist in two forms. A very small fraction of amino acid molecules will be neutral, with a deprotonated amino group and a protonated carboxylic acid group. However, the overwhelming majority of molecules will be in a "Zwitterion" tautomer, with a positive charge on the (protonated) amino group and a negative charge on the (deprotonated) carboxylate group. Amino acids are linked via a Peptide bond. (From an organic chemistry perspective, a peptide bond is a type of amide group.) A peptide bond consists of a carbonyl group's carbon atom directly bound to the nitrogen atom of a secondary amine. A peptide chain will have an unbound amino group free at one end (called the N-terminus) and a single free carboxylate group at the other end (called the C-terminus). The peptide bond is planar, because resonance between the carbonyl group and the amino nitrogen lends the C-N bond a partial double-bond character. (It is possible to draw a resonance structure with a double bond between the carbon atom and the nitrogen atom, with a formal negative charge on the oxygen atom and a formal positive charge on the nitrogen.) This prevents rotation around the C-N bond, locking the peptide bond in the "trans" conformation, and holding six atoms in a plane: the α-carbon of one amino acid, the carbonyl carbon and oxygen atoms, the amino nitrogen and hydrogen atoms, and the &alpha-carbon of the second amino acid are all co-planar. The structure of a peptide chain's "backbone" can be described uniquely by the torsion angles between adjacent peptide units. The φ-angle is the torsion angle between the α-carbon and the carbonyl carbon of one amino acid; the ψ-angle is the torsion angle between the amino nitrogen and the α-carbon. It is important to note that the formation of this peptide bond is highly unfavorable under standard conditions. However, in the human body, there are enzymes that assist in facilitating this reaction, making peptide bond formation and proteins possible. Chemical classification of aminoacids. The 20 amino acids encoded by the genetic code are:<br> They are not shown in their zwitterionic state for clarity. At physiological pH (~pH 6.8), all would be so, except proline, which is a five-membered ring. Charged side-chains are shown ionic when they exist as such at physiological pH. Protein Electrophoresis. Protein Electrophoresis is a method in which a mixture of proteins can be separated and analyzed. Electrophoresis is based on the mobility of ions in an electric field. The charge distribution of the molecules is critical in the separation of all electrophoresis. In an electric field, electrophoresis is a passage of charged molecules in solution. Positively charged ions have tendency toward a negative electrode and inversely, negatively charged ions have tendency toward a positive electrode. The molecular weight results to a molecular friction which is directly proportional to the molecular charge and its voltage and inversely proportional to a molecule's mobility in an electric field. Gel electrophoresis is performed to analyze the molecular weights and the charge of the protein and is mostly used in electrophoresis of the protein. The gel electrophoresis is carried out in a thin piece of polyacrylamide. The crosslinked of acrylamide and N,N'-methylene-bis-acrylamide forms the polyacrylamide by polymerization. The size has an important effect in the movement of the protein molecule. The smaller molecules of protein would result in a faster passage of the molecules through the gel pores. The separation of the protein molecules in the gel affects the protein activity. In this process, first the protein reduces the disruption of disulfide bonding by heating which results in purification and denaturalization. Next, the sodium dodecyl sulfate, abbreviated as SDS, (and ionic detergent) is added. SDS is an anionic detergent which dissolve hydrophobic molecules and denatures protein molecules without breaking peptide bonds. This result in the dislocation of the structure of the protein changes the secondary, tertiary and quaternary to the primary structure with negative charge. Then the protein is passed through the gel. For denatured proteins, SDS can form a steady charge mass ratio in binding with proteins. The polyacrylamide gel electrophoresis is a very sensitive method capable of a bearing a high resolution and it is analytically used in the studies in the separation techniques. Isoelectric point. The isoelectric point (pI) is the pH-value in which a protein is neutral, that is, has zero net charge. To be clear, it is not the pH value where a protein has all bases deprotonated and all acids protonated, but rather the value where positive and negative charges cancel out to zero. Calculating pI: An aminoacid with n ionizable groups with their respective pKa values pK1, pK2, ... pkn will have the pI equal to the average of the group pkas: pI=(pK1+pK2+...+pkn)/n Most proteins have many ionizable sidechains in addition to their amino- and carboxy- terminal groups. The pI is different for each protein and it can be theoretically calculated according Henderson-Hasselbalch equation if we know amino acids composition of protein. In order to experimentally determine a protein's pI 2-Dimensional Electrophoresis (2-DE) can be used. The proteins of a cell lysate are applied to a pH immobilized gradient strip, upon electrophoresis the proteins migrate to their pI within the strip. The second dimension of 2-DE is the separation of proteins by MW using a SDS-gel. To clearly understand isoelectric point, you have to keep in mind that the positively charged groups are balanced by the negatively charged groups. For simple amino acid " alanine", the isoelectric point is an average of the PKa of the carboxyl with PK1 of (2.34) and ammonium with PK2 of (9.69) groups. So, the PI for the simple amino acid " alanine" is calculated as:(2.34+9.69)/2 which is equal to 6.02. When additional basic or acidic groups are added as a side-chain functions, the isoelectric point pI will be the average of the pKa's of the most similar acids. Example for this concept could be aspartic acid in which the similar acids are alpha-carboxyl function with pKa of 2.1 and the side chain carboxyl function with a pKa of 3.9. Thus, the pI for aspartic acid is (2.1+3.9)/2=3.0. Another example is arginine, its similar acids are guanidinium on the side chain with pKa of 12.5 and the alpha-ammonium function with pKa of 9.0. Thus the calculated pI for arginine= (12.5+9.0)/2=10.75. pI does not has a unite. The peptidic bond. Two amino acid molecules can be covalently joined through a substituted amide linkage, termed a peptide bond, to yield a dipeptide. this link is formed by dehydration (removal of the water molecules - one hydrogen atom from one amino acid and an OH group from the other). This process can then continue to join other amino acids and yield in an amino acid chain. When there are few amino acids in a chain, it is called an oligopeptide, when there are many it is called a polypeptide. although the terms "protein" and "polypeptide" are sometimes used to describe the same thing, the term polypeptide is generally used when the molecular weight of the chain is below 10,000. An amino acid unit in a peptide is often called a residue. Disulfide bonds. Disulfide bonds form between the sulfur atoms of two cysteine side chains in a protein. The side chains undergo a reversible oxidation of the sulfhydryl groups of the cysteine, which results in covalent bonding of the sulfur atoms (S-S). This bonding is called a "disulfide bridge". Typically, disulfides don't form on the surface of proteins because of the presence of reducing agents in the cytoplasm. These bonds are of great importance concerning the shaping of protein structure; their formation guides the folding of peptide chains as the proteins are produced. Structural proteins that have to be rigorously stable (for example, Keratin, which is found in nail, horn and crustacean shell) often contain a large number of disulfide bonds. Post-translational modifications. After a protein is synthesized inside the cell, it is usually modified by the addition of extra functional groups to the polypeptide chain. These can be sugar or phosphate groups and may confer to the protein special functions such as: the ability to recognize other molecules, to integrate in the plasma membrane, to catalyze biochemical reactions, and various other processes. It is in the interest of the biochemist to understand what proteins are modified, what the modification is, and where it is located. An easy way to do this is by using mass spectrometry. In a sample of protein submitted to mass spectrometry you will see both modified and unmodified protein signals. The change in mass between these signals will correspond to the change in mass of your protein due to your post-translation modification.

Biochemistry/Proteins/Protein structure and folding. What is protein folding? Protein folding is commonly a fast or very fast process, often but not always reversible, taking no more than a few milliseconds to occur. It can be viewed as a complex compromise between the different chemical interactions that can happen between the amino acid sidechains, the amidic backbone and the solvent. There are literally millions of possible three-dimensional configurations, often with minimal energetic differences between them. That's why we're still almost unable to predict the folding of a given polypeptide chain "ab initio".Protein folding problem is that scientist still has failed to crack the code that governs folding. Moreover, the ability of biological polymers such as protein fold into well-defined structures is remarkable thermodynamically. An unfolded polymer exists are random coils, each copy of an unfolded polymers will have different conformation, yielding a mixture of many possible conformation. Folding depends only on primary structure. In 1954 Christian Anfinsen demonstrated that the folding of a protein in a given environment depends only on its primary structure - the amino acid sequence. This conclusion was by no means obvious, given the complexity of the folding process and the paucity of biochemical knowledge at the time. The process of folding often begins co-translationally, so that the N-terminus of the protein begins to fold while the C-terminal portion of the protein is still being synthesized by the ribosome. Specialized proteins called chaperones assist in the folding of other proteins. Meanwhile, protein folding is a thermodynamically driven process: that is, proteins fold by reaching their thermodynamically most stable structure. The path followed by the protein in the potential energy landscape is far from obvious, however. Many local and non-local interactions take part in the process, and the space of possible structures is enormous. As of today molecular dynamics simulations are giving invaluable hints on the first stages of the folding process. It is known now that the unfolded state still retains key long-range interactions and that the local propensity of the sequence to fold in a given secondary structure element narrow the "search" in the so-called conformational space. This seems to mean that biological proteins somehow evolved to properly fold. In fact, many random amino acid sequences only acquire ill-defined structures (molten globules) or no structure at all. There are some general rules, however. Hydrophobic amino acids will tend to be kept inside the structure, with little or no contact with the surrounding water; conversely, polar or charged amino acids will be often exposed to solvent. Very long proteins will often fold in various distinct modules, instead of in a single large structure. The Ramachandran plot. The Ramachandran plot was invented by professor G.N.Ramachandran, a very eminent scientist from India. He discovered the triple helix structure of collagen in 1954 a year after the double helix structure was put forth. In a polypeptide the main chain N-Calpha and Calpha-C bonds relatively are free to rotate. These rotations are represented by the torsion angles phi and psi, respectively. G N Ramachandran used computer models of small polypeptides to systematically vary phi and psi with the objective of finding stable conformations. For each conformation, the structure was examined for close contacts between atoms. Atoms were treated as hard spheres with dimensions corresponding to their van der Waals radii. Therefore, phi and psi angles which cause spheres to collide correspond to sterically disallowed conformations of the polypeptide backbone Intramolecular forces in protein folding. The tertiary structure is held together by hydrogen bonds, hydrophilic interactions, ionic interactions, and/or disulfide bonds. The protein folding problem. The protein folding problem relates to what is known as the Levinthal paradox. Levinthal calculated that if a fairly small protein is composed of 100 amino acids and each amino acid residue has only 3 possible conformations (an underestimate) then the entire protein can fold into 3100-1 or 5x1047 possible conformations. Even if it takes only 10-13 of a second to try each conformation it would take 1027 years to try them all. Obviously a protein doesn't take that long to fold, so randomly trying out all possible conformations is not the way proteins fold. Since most proteins fold on a timescale of the order of milliseconds it is clear that the process is directed in some manner dependent on the constituents of the chain. The protein folding problem which has perplexed scientists for over thirty years is that of understanding how the tertiary structure of a protein is related to its primary structure, because it has been proven that the primary structure of a protein holds the only information necessary for the protein to fold. Ultimately the aim is also to be able to predict what pathway the protein will take. Folding in extreme environments. Most proteins are not capable of maintaining their three-dimensional shape when they are exposed to environmental extremes such as a low or high pH, or a highly variable temperature. Changes in the pH of the proteins environment may alter the charges on the amino acid side chains that form the whole protein, so that repulsive or attractive forces may form, altering the secondary and tertiary structure of the protein as a whole, as a result the shape of the enzyme is warped, and the now non-functional protein is said to be denatured. A high or exceptionally low temperature can cause the constitutive bonds on the protein to be broken, again resulting in the protein being rendered non-functional. It is important to note that certain proteins, mainly digestive enzymes such as trypsin, are capable of with-standing a pH as low as 1. If the pH of such an enzymes environment were to increase to approximately pH5, it would be inactivated. Protein misfolding. The way that a protein folds is one of the most important factors influencing its properties, as this is what determines which active groups are exposed for interaction. If the protein misfolds, its properties can be markedly changed. One example of this is in Transmissible Spongiform Encephalopathies, such as BSE, and Scrapie. In these, the prion protein, which is involved in the brain's copper metabolism, misfolds, and starts forming plaques, which destroy brain tissue. =Protein structural levels= Biochemists refer to four distinct aspects of a protein's structure: Primary structure. Primary structure is practically a synonym of the amino acid sequence. It can also contain informations on amino acids linked by peptide bonds. Primary structure is typically written as a string of three letter sequences, each representing an amino acid. Peptides and proteins must have the correct sequence of amino acids. Secondary structure. "Secondary structure" elements are elementary structural patterns that are present in most,if not all,known proteins. These are highly patterned sub-structures --alpha helix and beta sheet-- consisting of loops between elements or segments of polypeptide chain that assume no stable shape. Secondary structure elements, when mapped on the sequence and depicted in the relative position they have in respect to each other, define the "topology" of the protein. It is also relevant to note that hydrogen bonding between residues is the cause for secondary structure features; secondary structure is usually described to beginning biochemists as (almost) entirely independent of residue side-chain interactions. Tertiary structure. "Tertiary structure" is the name given to refer to the overall shape of a single protein molecule. Although tertiary structure is sometimes described (especially to beginning biology and biochemistry students) as being a result of interactions between amino acid residue side chains, a more correct understanding of tertiary structure is the interactions between elements of secondary protein structure, i.e. alpha-helices and beta-pleated sheets. Tertiary structure is often referred to as the "fold structure" of a protein, since it is the result of the complex three-dimensional interplay of other structural and environmental elements. Super-tertiary structure (protein modules). Some literature refers to elements of super-tertiary structure, which often refers to elements of folding that, for whatever reason, do not neatly fit into the category of tertiary structure. Often this level of distinction is saved for graduate level coursework. Protein denaturation can be a reversible or an irreversible process, i.e., it may be possible or impossible to make the protein regain its original spatial conformation. Quaternary structure. Quaternary structure is the shape or structure that results from the union of more than one protein molecule, usually called "subunit proteins" "subunits" in this context, which function as part of the larger assembly or protein complex. And it refers to the regular association of two or more polypeptide chains to form a complex. A multi-subunit protein may be composed of two or more identical polypeptides, or it may include different polypeptides. Quaternary structure tends to be stabilized mainly by weak interactions between residues exposed on surfaces polypeptides within a complex. =Secondary structure elements= The alpha helix. The alpha helix is a periodic structure formed when main-chain atoms from residues spaced four residues apart hydrogen bond with one another. This gives rise to a helical structure, which in natural proteins is always right-handed. Each turn of the helix comprises 3.6 amino acids. Alpha helices are stiff, rod-like structures which are found in many unrelated proteins. One feature of these structures is that they tend to show a bias in the distribution of hydrophobic residues such that they tend to occur primarily on one face of the helix. The amino acids in an α helix are arranged in a helical structure, about 5 Å wide. Each amino acid results in a 100° turn in the helix, and corresponds to a translation of 1.5Å along the helical axis. The helix is tightly packed; there is almost no free space within the helix. All amino acid side-chains are arranged at the outside of the helix. The N-H group of amino acid (n) can establish a hydrogen bond with the C=O group of amino acid (n+4). Short polypeptides usually are not able to adopt the alpha helical structure, since the entropic cost associated with the folding of the polypeptide chain is too high. Some amino acids (called "helix breakers") like proline will disrupt the helical structure. Ordinarily, a helix has a buildup of positive charge at the N-terminal end and negative charge at the C-terminal end which is a destabilizing influence. As a result, α helices are often capped at the N-terminal end by a negatively charged amino acid (like glutamic acid) in order to stabilize the helix dipole. Less common (and less effective) is C-terminal capping with a positively charged protein like lysine. α helices have particular significance in DNA binding motifs, including helix-turn-helix motifs, leucine zipper motifs and zinc finger motifs. This is because of the structural coincidence of the α helix diameter of 12Å being the same as the width of the major groove in B-form DNA. α helices are one of the basic structural elements in proteins, together with beta sheets. The peptide backbone of an α helix has 3.6 amino acids per turn. The beta sheet. The β sheet (also β-pleated sheet) is a commonly occurring form of regular secondary structure in proteins, first proposed by Linus Pauling and Robert Corey in 1951. It consists of two or more amino acid sequences within the same protein that are arranged adjacently and in parallel, but with alternating orientation such that hydrogen bonds can form between the two strands. The amino acid chain is almost fully extended throughout a β strand, lessening the probability of bulky steric clashes. The Ramachandran plot shows optimal conformation with angle phi = -120 to -60 degrees and angle psi = 120 to 160 degrees. The N-H groups in the backbone of one strand establish hydrogen bonds with the C=O groups in the backbone of the adjacent, parallel strand(s). The cumulative effect of multiple such hydrogen bonds arranged in this way contributes to the sheet's stability and structural rigidity and integrity: for example, cellulose's beta-1,4 glucose structure. The side chains from the amino acid residues found in a β sheet structure may also be arranged such that many of the adjacent side chains on one side of the sheet are hydrophobic, while many of those adjacent to each other on the alternate side of the sheet are polar or charged (hydrophilic). Some sequences involved in a β sheet, when traced along the backbone, take a "hairpin turn" in orientation (direction), sometimes through one or more prolines. The α-C atoms of adjacent strands stand 3.5Å apart. In addition to the parallel beta sheet, there is also the anti-parallel beta sheet. The hydrogen-bond still exists in this conformation but the main difference lies in the directionality of the protein. In this conformation the proteins run in opposite directions, but the resulting hydrogen bonds connect directly across from one another instead of the diagonal. In short, beta sheets can be purely parallel, anti-parallel, or even mixed. The random coil. A protein that completely lacks secondary structure is a random coil. In random coil, the only fixed relationship between amino acids is that between adjacent residues through the peptide bond. As a result, random coil can be detected from the absence of the signals in a multidimensional nuclear magnetic resonance experiment that depend on particular peptide-peptide interactions. Likewise in the images produced in crystallography experiments, pieces of random coil appear simply as an absence of "electron density" or contrast. Random coil is also easily distinguished by circular dichroism. Denaturing reduces a protein entirely to random coil. Less common secondary structure elements. Certain other periodic structures rarely appear in proteins, some of which are similar to the more common types. For example, two variants of the alpha-helix also occur, the 3-10 helix and the pi helix. These have helical pitches of 3 and 4.4 residues per turn respectively, corresponding to hydrogen bonds forming between residue i and i+3 for the 3-10 helix and between i and i+5 for the pi helix. Both are usually very short (1 turn or so) and have been seen only at the ends of alpha helices. =Tertiary structure elements= Alpha-only structures. Structures that contain only Alpha Helices Beta-only structures. =Common folds and modules= The Rossman fold. The Rossmann fold is a protein structural motif found in proteins that bind nucleotides, especially the cofactor NAD. The structure is composed of three or more parallel beta strands linked by two alpha helices in the topological order beta-alpha-beta-alpha-beta. Because each Rossmann fold can bind one nucleotide, binding domains for dinucleotides such as NAD consist of two paired Rossmann folds that each bind one nucleotide moiety of the cofactor molecule. Single Rossmann folds can bind mononucleotides such as the cofactor FMN. The motif is named for Michael Rossmann who first pointed out that this was a frequently occurring motif in nucleotide binding proteins. Keratin. Found in hair and in the palms of human hands.

Guitar/Picking and Plucking. There are two major methods of right hand (for right handed players) techniques namely, either by using a pick (also called a plectrum) or fingers. The plectrum is very common in rock, country and pop music, where it is considered convenient for strumming and louder guitar sound. Use of fingers is most common among classical guitarists and flamenco players, as combination of strings better executed using the right hand fingers, and generally have softer sound than the pick. Other than classical guitarists and flamenco players, use of a pick or fingers is a matter of personal preference. Striking. Using a pick. The primary advantages of the pick are its speed, its ease of striking large chords and, because the fingernails and fingertips are not involved, its preservation of player's picking hand. Furthermore, use of a pick makes a louder and brighter sound. Its primary disadvantage is its imprecision, making muting strings necessary. Also, if the player wishes to switch to the tapping style, he or she can tap with or with out the pick: to tap with the pick just put it on its side and tap it on the desired fret. However, tapping with a pick makes it harder to tap on multiple strings. Finger Strumming. Players wishing not to use a pick may try finger "strumming". This is accomplished by holding the picking hand's first finger to the thumb, much as one might hold a pick, and striking the strings with the first fingernail. Another way is to do all down strokes with the thumb and all upstrokes with the index finger; like one is 'petting' the strings. Apoyando Strikes. Apoyando, or splinter rested, involves the finger picking through a string such that the finger stops when resting on the next string. This technique produces a strong, loud tone, and is considered the opposite of Tirando. Tirando Strikes. When performing a tirando, or shooting splinter strike, the finger does not affect the next string at all. This is the opposite of apoyando. Fingerpicking. Fingerpicking is a method of playing the guitar where you use your thumb and at least one other finger to pick or pluck notes, using your fingernails, fingerpicks or fingertips. Talented players can use all five fingers on their picking hand, but many players only use four fingers and use their pinky finger as a brace on the guitar. Most classical guitarists alter the shape of their picking hand fingernails for the purpose of producing a desired sound, however this is not necessary in non-classical music; one can purchase fingerpicks to fit on the hand. Generally fingerpicking involves picking through chords organized in a melody. Fingerpicking is used extensively in folk guitar and classical guitar, but it is also common in other genres. Fingerpicking is surprisingly easy on an electric guitar, which is strange because fingerpicking is often regarded as an acoustic style. The player may hold his or her picking hand's fourth finger against the right edge (left edge on a left-handed guitar), and if it is held straight and steady, this technique may be used to brace the hand. This technique is called anchoring, and is frowned upon by some players. It is possible on acoustic guitars by using the bridge similarly, but this is not as effective as it will deaden the sound. Classical guitarists never anchor while playing. When strumming with individual fingers, general rule is move the wrist only if the thumb is used, while if any other finger is used, only said finger will be used. When you start trying to learn, your finger coordination will be terrible and it is easy to be discouraged. It takes several weeks to let your muscles develop, but if you practice using all your fingers at once your overall dexterity will increase much faster. Classical picking. In classical guitar repertoire, there will be a "PIMA" marking for the picking hand fingers (right hand for right handed players), which indicate which finger to use: These four are the ones that are used most frequently. Sometimes, the fourth finger is used, in which it is marked either C, X or E. Typically, the thumb has a down-picking motion and the fingers have an up-picking motion. Clawhammer and frailing. Clawhammer, sometimes known as frailing, is a method generally used with the five-string banjo and is characteristic of traditional Appalachian folk music of the U.S. It is primarily a down-picking style, and the hand assumes a claw-like shape and the strumming finger is kept fairly stiff, striking the strings by the motion of the hand at the wrist and elbow, rather than a flicking motion by the finger. Typically, only the thumb and second or first finger are used and the finger always downpicks, flicking the string with the back of the fingernail. A common characteristic of clawhammer patterns is the thumb does not pick on the downbeat, as one might in typical finger-picking patterns for guitar. For example, this is a common, basic time signature|2/4 pattern: Here, the thumb plays the high drone (fifth string) on the second "and" of "one and two "and"". This combined with the second finger strumming provides a characteristic "bum-ditty bum-ditty" sound. Some people, however, make a distinction between frailing and clawhammer: Travis Picking. Another well known style of finger picking is called Travis picking, named after Merle Travis who was a country singer known for his legendary picking skills. When picking, you use your thumb and first finger to hit notes at the same time, creating a double stop or interval, and then continue picking with the first finger. Usually the thumb is responsible for picking the bass line, while the first/second finger is for melody. Skilled players can carry two separate melodies with the upper and lower strings. You can create impressive rhythms playing with just your thumb and first finger, but to really become talented you must practice using more fingers. For example, Chet Atkins expanded to use all three fingers, with thumb for bass line. Rasgueado. The rasgueado or splinter striking technique originated from Spanish flamenco music, and usually refers to three or four fingers and sometimes the thumb striking the strings in quick succession. The notes quickly follow one another and produce a "rattling" or cascading effect. Scruggs style. Scruggs-style finger-picking is a syncopated, five-string banjo style used in bluegrass music. It is played with thumb, first and second fingers; the fourth and/or third fingers are typically braced against the head of the instrument. The strings are picked rapidly in repetitive sequences or rolls; the same string is not typically picked twice in succession. Melody notes are interspersed among arpeggios, and musical phrases typically contain long series of staccato notes, often played at very rapid tempos. The music is generally syncopated, and may have a subtle swing or shuffle feel, especially on mid-tempo numbers. The result is lively, rapid music, which lends itself both as an accompaniment to other instruments and as a solo. Scruggs style picking was popularized by Earl Scruggs in the early 1940's in rural North Carolina. Tapping. Tapping is a style of playing where notes are created by quickly pressing, or tapping, the string down on the fret that you want to play. Usually tapping involves both hands, and most often it is on an electric guitar. It is possible to tap on an acoustic, but you cannot hear the notes as clearly as on an electric.