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# Copyright 2023 The HuggingFace Team. All rights reserved. | |
# | |
# Licensed under the Apache License, Version 2.0 (the "License"); | |
# you may not use this file except in compliance with the License. | |
# You may obtain a copy of the License at | |
# | |
# http://www.apache.org/licenses/LICENSE-2.0 | |
# | |
# Unless required by applicable law or agreed to in writing, software | |
# distributed under the License is distributed on an "AS IS" BASIS, | |
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. | |
# See the License for the specific language governing permissions and | |
# limitations under the License. | |
import inspect | |
from typing import Any, Callable, Dict, List, Optional, Tuple, Union | |
import numpy as np | |
import PIL.Image | |
import torch | |
from transformers import ( | |
CLIPImageProcessor, | |
CLIPTextModel, | |
CLIPTextModelWithProjection, | |
CLIPTokenizer, | |
CLIPVisionModelWithProjection, | |
) | |
from diffusers.image_processor import PipelineImageInput, VaeImageProcessor | |
from diffusers.loaders import ( | |
FromSingleFileMixin, | |
IPAdapterMixin, | |
StableDiffusionXLLoraLoaderMixin, | |
TextualInversionLoaderMixin, | |
) | |
from diffusers.models import AutoencoderKL, ImageProjection, UNet2DConditionModel | |
from diffusers.models.attention_processor import ( | |
AttnProcessor2_0, | |
FusedAttnProcessor2_0, | |
LoRAAttnProcessor2_0, | |
LoRAXFormersAttnProcessor, | |
XFormersAttnProcessor, | |
) | |
from diffusers.models.lora import adjust_lora_scale_text_encoder | |
from diffusers.schedulers import KarrasDiffusionSchedulers | |
from diffusers.utils import ( | |
USE_PEFT_BACKEND, | |
deprecate, | |
is_invisible_watermark_available, | |
is_torch_xla_available, | |
logging, | |
replace_example_docstring, | |
scale_lora_layers, | |
unscale_lora_layers, | |
) | |
from diffusers.utils.torch_utils import randn_tensor | |
from diffusers.pipelines.pipeline_utils import DiffusionPipeline | |
if is_torch_xla_available(): | |
import torch_xla.core.xla_model as xm | |
XLA_AVAILABLE = True | |
else: | |
XLA_AVAILABLE = False | |
logger = logging.get_logger(__name__) # pylint: disable=invalid-name | |
EXAMPLE_DOC_STRING = """ | |
Examples: | |
```py | |
>>> import torch | |
>>> from diffusers import StableDiffusionXLInpaintPipeline | |
>>> from diffusers.utils import load_image | |
>>> pipe = StableDiffusionXLInpaintPipeline.from_pretrained( | |
... "stabilityai/stable-diffusion-xl-base-1.0", | |
... torch_dtype=torch.float16, | |
... variant="fp16", | |
... use_safetensors=True, | |
... ) | |
>>> pipe.to("cuda") | |
>>> img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" | |
>>> mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" | |
>>> init_image = load_image(img_url).convert("RGB") | |
>>> mask_image = load_image(mask_url).convert("RGB") | |
>>> prompt = "A majestic tiger sitting on a bench" | |
>>> image = pipe( | |
... prompt=prompt, image=init_image, mask_image=mask_image, num_inference_steps=50, strength=0.80 | |
... ).images[0] | |
``` | |
""" | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.rescale_noise_cfg | |
def rescale_noise_cfg(noise_cfg, noise_pred_text, guidance_rescale=0.0): | |
""" | |
Rescale `noise_cfg` according to `guidance_rescale`. Based on findings of [Common Diffusion Noise Schedules and | |
Sample Steps are Flawed](https://arxiv.org/pdf/2305.08891.pdf). See Section 3.4 | |
""" | |
std_text = noise_pred_text.std(dim=list(range(1, noise_pred_text.ndim)), keepdim=True) | |
std_cfg = noise_cfg.std(dim=list(range(1, noise_cfg.ndim)), keepdim=True) | |
# rescale the results from guidance (fixes overexposure) | |
noise_pred_rescaled = noise_cfg * (std_text / std_cfg) | |
# mix with the original results from guidance by factor guidance_rescale to avoid "plain looking" images | |
noise_cfg = guidance_rescale * noise_pred_rescaled + (1 - guidance_rescale) * noise_cfg | |
return noise_cfg | |
def mask_pil_to_torch(mask, height, width): | |
# preprocess mask | |
if isinstance(mask, (PIL.Image.Image, np.ndarray)): | |
mask = [mask] | |
if isinstance(mask, list) and isinstance(mask[0], PIL.Image.Image): | |
mask = [i.resize((width, height), resample=PIL.Image.LANCZOS) for i in mask] | |
mask = np.concatenate([np.array(m.convert("L"))[None, None, :] for m in mask], axis=0) | |
mask = mask.astype(np.float32) / 255.0 | |
elif isinstance(mask, list) and isinstance(mask[0], np.ndarray): | |
mask = np.concatenate([m[None, None, :] for m in mask], axis=0) | |
mask = torch.from_numpy(mask) | |
return mask | |
def prepare_mask_and_masked_image(image, mask, height, width, return_image: bool = False): | |
""" | |
Prepares a pair (image, mask) to be consumed by the Stable Diffusion pipeline. This means that those inputs will be | |
converted to ``torch.Tensor`` with shapes ``batch x channels x height x width`` where ``channels`` is ``3`` for the | |
``image`` and ``1`` for the ``mask``. | |
The ``image`` will be converted to ``torch.float32`` and normalized to be in ``[-1, 1]``. The ``mask`` will be | |
binarized (``mask > 0.5``) and cast to ``torch.float32`` too. | |
Args: | |
image (Union[np.array, PIL.Image, torch.Tensor]): The image to inpaint. | |
It can be a ``PIL.Image``, or a ``height x width x 3`` ``np.array`` or a ``channels x height x width`` | |
``torch.Tensor`` or a ``batch x channels x height x width`` ``torch.Tensor``. | |
mask (_type_): The mask to apply to the image, i.e. regions to inpaint. | |
It can be a ``PIL.Image``, or a ``height x width`` ``np.array`` or a ``1 x height x width`` | |
``torch.Tensor`` or a ``batch x 1 x height x width`` ``torch.Tensor``. | |
Raises: | |
ValueError: ``torch.Tensor`` images should be in the ``[-1, 1]`` range. ValueError: ``torch.Tensor`` mask | |
should be in the ``[0, 1]`` range. ValueError: ``mask`` and ``image`` should have the same spatial dimensions. | |
TypeError: ``mask`` is a ``torch.Tensor`` but ``image`` is not | |
(ot the other way around). | |
Returns: | |
tuple[torch.Tensor]: The pair (mask, masked_image) as ``torch.Tensor`` with 4 | |
dimensions: ``batch x channels x height x width``. | |
""" | |
# checkpoint. TOD(Yiyi) - need to clean this up later | |
deprecation_message = "The prepare_mask_and_masked_image method is deprecated and will be removed in a future version. Please use VaeImageProcessor.preprocess instead" | |
deprecate( | |
"prepare_mask_and_masked_image", | |
"0.30.0", | |
deprecation_message, | |
) | |
if image is None: | |
raise ValueError("`image` input cannot be undefined.") | |
if mask is None: | |
raise ValueError("`mask_image` input cannot be undefined.") | |
if isinstance(image, torch.Tensor): | |
if not isinstance(mask, torch.Tensor): | |
mask = mask_pil_to_torch(mask, height, width) | |
if image.ndim == 3: | |
image = image.unsqueeze(0) | |
# Batch and add channel dim for single mask | |
if mask.ndim == 2: | |
mask = mask.unsqueeze(0).unsqueeze(0) | |
# Batch single mask or add channel dim | |
if mask.ndim == 3: | |
# Single batched mask, no channel dim or single mask not batched but channel dim | |
if mask.shape[0] == 1: | |
mask = mask.unsqueeze(0) | |
# Batched masks no channel dim | |
else: | |
mask = mask.unsqueeze(1) | |
assert image.ndim == 4 and mask.ndim == 4, "Image and Mask must have 4 dimensions" | |
# assert image.shape[-2:] == mask.shape[-2:], "Image and Mask must have the same spatial dimensions" | |
assert image.shape[0] == mask.shape[0], "Image and Mask must have the same batch size" | |
# Check image is in [-1, 1] | |
# if image.min() < -1 or image.max() > 1: | |
# raise ValueError("Image should be in [-1, 1] range") | |
# Check mask is in [0, 1] | |
if mask.min() < 0 or mask.max() > 1: | |
raise ValueError("Mask should be in [0, 1] range") | |
# Binarize mask | |
mask[mask < 0.5] = 0 | |
mask[mask >= 0.5] = 1 | |
# Image as float32 | |
image = image.to(dtype=torch.float32) | |
elif isinstance(mask, torch.Tensor): | |
raise TypeError(f"`mask` is a torch.Tensor but `image` (type: {type(image)} is not") | |
else: | |
# preprocess image | |
if isinstance(image, (PIL.Image.Image, np.ndarray)): | |
image = [image] | |
if isinstance(image, list) and isinstance(image[0], PIL.Image.Image): | |
# resize all images w.r.t passed height an width | |
image = [i.resize((width, height), resample=PIL.Image.LANCZOS) for i in image] | |
image = [np.array(i.convert("RGB"))[None, :] for i in image] | |
image = np.concatenate(image, axis=0) | |
elif isinstance(image, list) and isinstance(image[0], np.ndarray): | |
image = np.concatenate([i[None, :] for i in image], axis=0) | |
image = image.transpose(0, 3, 1, 2) | |
image = torch.from_numpy(image).to(dtype=torch.float32) / 127.5 - 1.0 | |
mask = mask_pil_to_torch(mask, height, width) | |
mask[mask < 0.5] = 0 | |
mask[mask >= 0.5] = 1 | |
if image.shape[1] == 4: | |
# images are in latent space and thus can't | |
# be masked set masked_image to None | |
# we assume that the checkpoint is not an inpainting | |
# checkpoint. TOD(Yiyi) - need to clean this up later | |
masked_image = None | |
else: | |
masked_image = image * (mask < 0.5) | |
# n.b. ensure backwards compatibility as old function does not return image | |
if return_image: | |
return mask, masked_image, image | |
return mask, masked_image | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_img2img.retrieve_latents | |
def retrieve_latents( | |
encoder_output: torch.Tensor, generator: Optional[torch.Generator] = None, sample_mode: str = "sample" | |
): | |
if hasattr(encoder_output, "latent_dist") and sample_mode == "sample": | |
return encoder_output.latent_dist.sample(generator) | |
elif hasattr(encoder_output, "latent_dist") and sample_mode == "argmax": | |
return encoder_output.latent_dist.mode() | |
elif hasattr(encoder_output, "latents"): | |
return encoder_output.latents | |
else: | |
raise AttributeError("Could not access latents of provided encoder_output") | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.retrieve_timesteps | |
def retrieve_timesteps( | |
scheduler, | |
num_inference_steps: Optional[int] = None, | |
device: Optional[Union[str, torch.device]] = None, | |
timesteps: Optional[List[int]] = None, | |
**kwargs, | |
): | |
""" | |
Calls the scheduler's `set_timesteps` method and retrieves timesteps from the scheduler after the call. Handles | |
custom timesteps. Any kwargs will be supplied to `scheduler.set_timesteps`. | |
Args: | |
scheduler (`SchedulerMixin`): | |
The scheduler to get timesteps from. | |
num_inference_steps (`int`): | |
The number of diffusion steps used when generating samples with a pre-trained model. If used, | |
`timesteps` must be `None`. | |
device (`str` or `torch.device`, *optional*): | |
The device to which the timesteps should be moved to. If `None`, the timesteps are not moved. | |
timesteps (`List[int]`, *optional*): | |
Custom timesteps used to support arbitrary spacing between timesteps. If `None`, then the default | |
timestep spacing strategy of the scheduler is used. If `timesteps` is passed, `num_inference_steps` | |
must be `None`. | |
Returns: | |
`Tuple[torch.Tensor, int]`: A tuple where the first element is the timestep schedule from the scheduler and the | |
second element is the number of inference steps. | |
""" | |
if timesteps is not None: | |
accepts_timesteps = "timesteps" in set(inspect.signature(scheduler.set_timesteps).parameters.keys()) | |
if not accepts_timesteps: | |
raise ValueError( | |
f"The current scheduler class {scheduler.__class__}'s `set_timesteps` does not support custom" | |
f" timestep schedules. Please check whether you are using the correct scheduler." | |
) | |
scheduler.set_timesteps(timesteps=timesteps, device=device, **kwargs) | |
timesteps = scheduler.timesteps | |
num_inference_steps = len(timesteps) | |
else: | |
scheduler.set_timesteps(num_inference_steps, device=device, **kwargs) | |
timesteps = scheduler.timesteps | |
return timesteps, num_inference_steps | |
class StableDiffusionXLInpaintPipeline( | |
DiffusionPipeline, | |
TextualInversionLoaderMixin, | |
StableDiffusionXLLoraLoaderMixin, | |
FromSingleFileMixin, | |
IPAdapterMixin, | |
): | |
r""" | |
Pipeline for text-to-image generation using Stable Diffusion XL. | |
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the | |
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.) | |
The pipeline also inherits the following loading methods: | |
- [`~loaders.TextualInversionLoaderMixin.load_textual_inversion`] for loading textual inversion embeddings | |
- [`~loaders.FromSingleFileMixin.from_single_file`] for loading `.ckpt` files | |
- [`~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights`] for loading LoRA weights | |
- [`~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights`] for saving LoRA weights | |
- [`~loaders.IPAdapterMixin.load_ip_adapter`] for loading IP Adapters | |
Args: | |
vae ([`AutoencoderKL`]): | |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations. | |
text_encoder ([`CLIPTextModel`]): | |
Frozen text-encoder. Stable Diffusion XL uses the text portion of | |
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically | |
the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant. | |
text_encoder_2 ([` CLIPTextModelWithProjection`]): | |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of | |
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModelWithProjection), | |
specifically the | |
[laion/CLIP-ViT-bigG-14-laion2B-39B-b160k](https://huggingface.co/laion/CLIP-ViT-bigG-14-laion2B-39B-b160k) | |
variant. | |
tokenizer (`CLIPTokenizer`): | |
Tokenizer of class | |
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer). | |
tokenizer_2 (`CLIPTokenizer`): | |
Second Tokenizer of class | |
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer). | |
unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents. | |
scheduler ([`SchedulerMixin`]): | |
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of | |
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`]. | |
requires_aesthetics_score (`bool`, *optional*, defaults to `"False"`): | |
Whether the `unet` requires a aesthetic_score condition to be passed during inference. Also see the config | |
of `stabilityai/stable-diffusion-xl-refiner-1-0`. | |
force_zeros_for_empty_prompt (`bool`, *optional*, defaults to `"True"`): | |
Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config of | |
`stabilityai/stable-diffusion-xl-base-1-0`. | |
add_watermarker (`bool`, *optional*): | |
Whether to use the [invisible_watermark library](https://github.com/ShieldMnt/invisible-watermark/) to | |
watermark output images. If not defined, it will default to True if the package is installed, otherwise no | |
watermarker will be used. | |
""" | |
model_cpu_offload_seq = "text_encoder->text_encoder_2->image_encoder->unet->vae" | |
_optional_components = [ | |
"tokenizer", | |
"tokenizer_2", | |
"text_encoder", | |
"text_encoder_2", | |
"image_encoder", | |
"feature_extractor", | |
] | |
_callback_tensor_inputs = [ | |
"latents", | |
"prompt_embeds", | |
"negative_prompt_embeds", | |
"add_text_embeds", | |
"add_time_ids", | |
"negative_pooled_prompt_embeds", | |
"add_neg_time_ids", | |
"mask", | |
"masked_image_latents", | |
] | |
def __init__( | |
self, | |
vae: AutoencoderKL, | |
text_encoder: CLIPTextModel, | |
text_encoder_2: CLIPTextModelWithProjection, | |
tokenizer: CLIPTokenizer, | |
tokenizer_2: CLIPTokenizer, | |
unet: UNet2DConditionModel, | |
scheduler: KarrasDiffusionSchedulers, | |
image_encoder: CLIPVisionModelWithProjection = None, | |
feature_extractor: CLIPImageProcessor = None, | |
requires_aesthetics_score: bool = False, | |
force_zeros_for_empty_prompt: bool = True, | |
): | |
super().__init__() | |
self.register_modules( | |
vae=vae, | |
text_encoder=text_encoder, | |
text_encoder_2=text_encoder_2, | |
tokenizer=tokenizer, | |
tokenizer_2=tokenizer_2, | |
unet=unet, | |
image_encoder=image_encoder, | |
feature_extractor=feature_extractor, | |
scheduler=scheduler, | |
) | |
self.register_to_config(force_zeros_for_empty_prompt=force_zeros_for_empty_prompt) | |
self.register_to_config(requires_aesthetics_score=requires_aesthetics_score) | |
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1) | |
self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor) | |
self.mask_processor = VaeImageProcessor( | |
vae_scale_factor=self.vae_scale_factor, do_normalize=False, do_binarize=True, do_convert_grayscale=True | |
) | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_slicing | |
def enable_vae_slicing(self): | |
r""" | |
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to | |
compute decoding in several steps. This is useful to save some memory and allow larger batch sizes. | |
""" | |
self.vae.enable_slicing() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_slicing | |
def disable_vae_slicing(self): | |
r""" | |
Disable sliced VAE decoding. If `enable_vae_slicing` was previously enabled, this method will go back to | |
computing decoding in one step. | |
""" | |
self.vae.disable_slicing() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_tiling | |
def enable_vae_tiling(self): | |
r""" | |
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to | |
compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow | |
processing larger images. | |
""" | |
self.vae.enable_tiling() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_tiling | |
def disable_vae_tiling(self): | |
r""" | |
Disable tiled VAE decoding. If `enable_vae_tiling` was previously enabled, this method will go back to | |
computing decoding in one step. | |
""" | |
self.vae.disable_tiling() | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.encode_image | |
def encode_image(self, image, device, num_images_per_prompt, output_hidden_states=None): | |
dtype = next(self.image_encoder.parameters()).dtype | |
# print(image.shape) | |
if not isinstance(image, torch.Tensor): | |
image = self.feature_extractor(image, return_tensors="pt").pixel_values | |
image = image.to(device=device, dtype=dtype) | |
if output_hidden_states: | |
image_enc_hidden_states = self.image_encoder(image, output_hidden_states=True).hidden_states[-2] | |
image_enc_hidden_states = image_enc_hidden_states.repeat_interleave(num_images_per_prompt, dim=0) | |
uncond_image_enc_hidden_states = self.image_encoder( | |
torch.zeros_like(image), output_hidden_states=True | |
).hidden_states[-2] | |
uncond_image_enc_hidden_states = uncond_image_enc_hidden_states.repeat_interleave( | |
num_images_per_prompt, dim=0 | |
) | |
return image_enc_hidden_states, uncond_image_enc_hidden_states | |
else: | |
image_embeds = self.image_encoder(image).image_embeds | |
image_embeds = image_embeds.repeat_interleave(num_images_per_prompt, dim=0) | |
uncond_image_embeds = torch.zeros_like(image_embeds) | |
return image_embeds, uncond_image_embeds | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_ip_adapter_image_embeds | |
def prepare_ip_adapter_image_embeds(self, ip_adapter_image, device, num_images_per_prompt): | |
# if not isinstance(ip_adapter_image, list): | |
# ip_adapter_image = [ip_adapter_image] | |
# if len(ip_adapter_image) != len(self.unet.encoder_hid_proj.image_projection_layers): | |
# raise ValueError( | |
# f"`ip_adapter_image` must have same length as the number of IP Adapters. Got {len(ip_adapter_image)} images and {len(self.unet.encoder_hid_proj.image_projection_layers)} IP Adapters." | |
# ) | |
output_hidden_state = not isinstance(self.unet.encoder_hid_proj, ImageProjection) | |
# print(output_hidden_state) | |
image_embeds, negative_image_embeds = self.encode_image( | |
ip_adapter_image, device, 1, output_hidden_state | |
) | |
# print(single_image_embeds.shape) | |
# single_image_embeds = torch.stack([single_image_embeds] * num_images_per_prompt, dim=0) | |
# single_negative_image_embeds = torch.stack([single_negative_image_embeds] * num_images_per_prompt, dim=0) | |
# print(single_image_embeds.shape) | |
if self.do_classifier_free_guidance: | |
image_embeds = torch.cat([negative_image_embeds, image_embeds]) | |
image_embeds = image_embeds.to(device) | |
return image_embeds | |
# Copied from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl.StableDiffusionXLPipeline.encode_prompt | |
def encode_prompt( | |
self, | |
prompt: str, | |
prompt_2: Optional[str] = None, | |
device: Optional[torch.device] = None, | |
num_images_per_prompt: int = 1, | |
do_classifier_free_guidance: bool = True, | |
negative_prompt: Optional[str] = None, | |
negative_prompt_2: Optional[str] = None, | |
prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
lora_scale: Optional[float] = None, | |
clip_skip: Optional[int] = None, | |
): | |
r""" | |
Encodes the prompt into text encoder hidden states. | |
Args: | |
prompt (`str` or `List[str]`, *optional*): | |
prompt to be encoded | |
prompt_2 (`str` or `List[str]`, *optional*): | |
The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is | |
used in both text-encoders | |
device: (`torch.device`): | |
torch device | |
num_images_per_prompt (`int`): | |
number of images that should be generated per prompt | |
do_classifier_free_guidance (`bool`): | |
whether to use classifier free guidance or not | |
negative_prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
less than `1`). | |
negative_prompt_2 (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and | |
`text_encoder_2`. If not defined, `negative_prompt` is used in both text-encoders | |
prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
provided, text embeddings will be generated from `prompt` input argument. | |
negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
argument. | |
pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. | |
If not provided, pooled text embeddings will be generated from `prompt` input argument. | |
negative_pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` | |
input argument. | |
lora_scale (`float`, *optional*): | |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded. | |
clip_skip (`int`, *optional*): | |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that | |
the output of the pre-final layer will be used for computing the prompt embeddings. | |
""" | |
device = device or self._execution_device | |
# set lora scale so that monkey patched LoRA | |
# function of text encoder can correctly access it | |
if lora_scale is not None and isinstance(self, StableDiffusionXLLoraLoaderMixin): | |
self._lora_scale = lora_scale | |
# dynamically adjust the LoRA scale | |
if self.text_encoder is not None: | |
if not USE_PEFT_BACKEND: | |
adjust_lora_scale_text_encoder(self.text_encoder, lora_scale) | |
else: | |
scale_lora_layers(self.text_encoder, lora_scale) | |
if self.text_encoder_2 is not None: | |
if not USE_PEFT_BACKEND: | |
adjust_lora_scale_text_encoder(self.text_encoder_2, lora_scale) | |
else: | |
scale_lora_layers(self.text_encoder_2, lora_scale) | |
prompt = [prompt] if isinstance(prompt, str) else prompt | |
if prompt is not None: | |
batch_size = len(prompt) | |
else: | |
batch_size = prompt_embeds.shape[0] | |
# Define tokenizers and text encoders | |
tokenizers = [self.tokenizer, self.tokenizer_2] if self.tokenizer is not None else [self.tokenizer_2] | |
text_encoders = ( | |
[self.text_encoder, self.text_encoder_2] if self.text_encoder is not None else [self.text_encoder_2] | |
) | |
if prompt_embeds is None: | |
prompt_2 = prompt_2 or prompt | |
prompt_2 = [prompt_2] if isinstance(prompt_2, str) else prompt_2 | |
# textual inversion: procecss multi-vector tokens if necessary | |
prompt_embeds_list = [] | |
prompts = [prompt, prompt_2] | |
for prompt, tokenizer, text_encoder in zip(prompts, tokenizers, text_encoders): | |
if isinstance(self, TextualInversionLoaderMixin): | |
prompt = self.maybe_convert_prompt(prompt, tokenizer) | |
text_inputs = tokenizer( | |
prompt, | |
padding="max_length", | |
max_length=tokenizer.model_max_length, | |
truncation=True, | |
return_tensors="pt", | |
) | |
text_input_ids = text_inputs.input_ids | |
untruncated_ids = tokenizer(prompt, padding="longest", return_tensors="pt").input_ids | |
if untruncated_ids.shape[-1] >= text_input_ids.shape[-1] and not torch.equal( | |
text_input_ids, untruncated_ids | |
): | |
removed_text = tokenizer.batch_decode(untruncated_ids[:, tokenizer.model_max_length - 1 : -1]) | |
logger.warning( | |
"The following part of your input was truncated because CLIP can only handle sequences up to" | |
f" {tokenizer.model_max_length} tokens: {removed_text}" | |
) | |
prompt_embeds = text_encoder(text_input_ids.to(device), output_hidden_states=True) | |
# We are only ALWAYS interested in the pooled output of the final text encoder | |
pooled_prompt_embeds = prompt_embeds[0] | |
if clip_skip is None: | |
prompt_embeds = prompt_embeds.hidden_states[-2] | |
else: | |
# "2" because SDXL always indexes from the penultimate layer. | |
prompt_embeds = prompt_embeds.hidden_states[-(clip_skip + 2)] | |
prompt_embeds_list.append(prompt_embeds) | |
prompt_embeds = torch.concat(prompt_embeds_list, dim=-1) | |
# get unconditional embeddings for classifier free guidance | |
zero_out_negative_prompt = negative_prompt is None and self.config.force_zeros_for_empty_prompt | |
if do_classifier_free_guidance and negative_prompt_embeds is None and zero_out_negative_prompt: | |
negative_prompt_embeds = torch.zeros_like(prompt_embeds) | |
negative_pooled_prompt_embeds = torch.zeros_like(pooled_prompt_embeds) | |
elif do_classifier_free_guidance and negative_prompt_embeds is None: | |
negative_prompt = negative_prompt or "" | |
negative_prompt_2 = negative_prompt_2 or negative_prompt | |
# normalize str to list | |
negative_prompt = batch_size * [negative_prompt] if isinstance(negative_prompt, str) else negative_prompt | |
negative_prompt_2 = ( | |
batch_size * [negative_prompt_2] if isinstance(negative_prompt_2, str) else negative_prompt_2 | |
) | |
uncond_tokens: List[str] | |
if prompt is not None and type(prompt) is not type(negative_prompt): | |
raise TypeError( | |
f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !=" | |
f" {type(prompt)}." | |
) | |
elif batch_size != len(negative_prompt): | |
raise ValueError( | |
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:" | |
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches" | |
" the batch size of `prompt`." | |
) | |
else: | |
uncond_tokens = [negative_prompt, negative_prompt_2] | |
negative_prompt_embeds_list = [] | |
for negative_prompt, tokenizer, text_encoder in zip(uncond_tokens, tokenizers, text_encoders): | |
if isinstance(self, TextualInversionLoaderMixin): | |
negative_prompt = self.maybe_convert_prompt(negative_prompt, tokenizer) | |
max_length = prompt_embeds.shape[1] | |
uncond_input = tokenizer( | |
negative_prompt, | |
padding="max_length", | |
max_length=max_length, | |
truncation=True, | |
return_tensors="pt", | |
) | |
negative_prompt_embeds = text_encoder( | |
uncond_input.input_ids.to(device), | |
output_hidden_states=True, | |
) | |
# We are only ALWAYS interested in the pooled output of the final text encoder | |
negative_pooled_prompt_embeds = negative_prompt_embeds[0] | |
negative_prompt_embeds = negative_prompt_embeds.hidden_states[-2] | |
negative_prompt_embeds_list.append(negative_prompt_embeds) | |
negative_prompt_embeds = torch.concat(negative_prompt_embeds_list, dim=-1) | |
if self.text_encoder_2 is not None: | |
prompt_embeds = prompt_embeds.to(dtype=self.text_encoder_2.dtype, device=device) | |
else: | |
prompt_embeds = prompt_embeds.to(dtype=self.unet.dtype, device=device) | |
bs_embed, seq_len, _ = prompt_embeds.shape | |
# duplicate text embeddings for each generation per prompt, using mps friendly method | |
prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1) | |
prompt_embeds = prompt_embeds.view(bs_embed * num_images_per_prompt, seq_len, -1) | |
if do_classifier_free_guidance: | |
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method | |
seq_len = negative_prompt_embeds.shape[1] | |
if self.text_encoder_2 is not None: | |
negative_prompt_embeds = negative_prompt_embeds.to(dtype=self.text_encoder_2.dtype, device=device) | |
else: | |
negative_prompt_embeds = negative_prompt_embeds.to(dtype=self.unet.dtype, device=device) | |
negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1) | |
negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1) | |
pooled_prompt_embeds = pooled_prompt_embeds.repeat(1, num_images_per_prompt).view( | |
bs_embed * num_images_per_prompt, -1 | |
) | |
if do_classifier_free_guidance: | |
negative_pooled_prompt_embeds = negative_pooled_prompt_embeds.repeat(1, num_images_per_prompt).view( | |
bs_embed * num_images_per_prompt, -1 | |
) | |
if self.text_encoder is not None: | |
if isinstance(self, StableDiffusionXLLoraLoaderMixin) and USE_PEFT_BACKEND: | |
# Retrieve the original scale by scaling back the LoRA layers | |
unscale_lora_layers(self.text_encoder, lora_scale) | |
if self.text_encoder_2 is not None: | |
if isinstance(self, StableDiffusionXLLoraLoaderMixin) and USE_PEFT_BACKEND: | |
# Retrieve the original scale by scaling back the LoRA layers | |
unscale_lora_layers(self.text_encoder_2, lora_scale) | |
return prompt_embeds, negative_prompt_embeds, pooled_prompt_embeds, negative_pooled_prompt_embeds | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs | |
def prepare_extra_step_kwargs(self, generator, eta): | |
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature | |
# eta (Ξ·) is only used with the DDIMScheduler, it will be ignored for other schedulers. | |
# eta corresponds to Ξ· in DDIM paper: https://arxiv.org/abs/2010.02502 | |
# and should be between [0, 1] | |
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys()) | |
extra_step_kwargs = {} | |
if accepts_eta: | |
extra_step_kwargs["eta"] = eta | |
# check if the scheduler accepts generator | |
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys()) | |
if accepts_generator: | |
extra_step_kwargs["generator"] = generator | |
return extra_step_kwargs | |
def check_inputs( | |
self, | |
prompt, | |
prompt_2, | |
image, | |
mask_image, | |
height, | |
width, | |
strength, | |
callback_steps, | |
output_type, | |
negative_prompt=None, | |
negative_prompt_2=None, | |
prompt_embeds=None, | |
negative_prompt_embeds=None, | |
callback_on_step_end_tensor_inputs=None, | |
padding_mask_crop=None, | |
): | |
if strength < 0 or strength > 1: | |
raise ValueError(f"The value of strength should in [0.0, 1.0] but is {strength}") | |
if height % 8 != 0 or width % 8 != 0: | |
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.") | |
if callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0): | |
raise ValueError( | |
f"`callback_steps` has to be a positive integer but is {callback_steps} of type" | |
f" {type(callback_steps)}." | |
) | |
if callback_on_step_end_tensor_inputs is not None and not all( | |
k in self._callback_tensor_inputs for k in callback_on_step_end_tensor_inputs | |
): | |
raise ValueError( | |
f"`callback_on_step_end_tensor_inputs` has to be in {self._callback_tensor_inputs}, but found {[k for k in callback_on_step_end_tensor_inputs if k not in self._callback_tensor_inputs]}" | |
) | |
if prompt is not None and prompt_embeds is not None: | |
raise ValueError( | |
f"Cannot forward both `prompt`: {prompt} and `prompt_embeds`: {prompt_embeds}. Please make sure to" | |
" only forward one of the two." | |
) | |
elif prompt_2 is not None and prompt_embeds is not None: | |
raise ValueError( | |
f"Cannot forward both `prompt_2`: {prompt_2} and `prompt_embeds`: {prompt_embeds}. Please make sure to" | |
" only forward one of the two." | |
) | |
elif prompt is None and prompt_embeds is None: | |
raise ValueError( | |
"Provide either `prompt` or `prompt_embeds`. Cannot leave both `prompt` and `prompt_embeds` undefined." | |
) | |
elif prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)): | |
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}") | |
elif prompt_2 is not None and (not isinstance(prompt_2, str) and not isinstance(prompt_2, list)): | |
raise ValueError(f"`prompt_2` has to be of type `str` or `list` but is {type(prompt_2)}") | |
if negative_prompt is not None and negative_prompt_embeds is not None: | |
raise ValueError( | |
f"Cannot forward both `negative_prompt`: {negative_prompt} and `negative_prompt_embeds`:" | |
f" {negative_prompt_embeds}. Please make sure to only forward one of the two." | |
) | |
elif negative_prompt_2 is not None and negative_prompt_embeds is not None: | |
raise ValueError( | |
f"Cannot forward both `negative_prompt_2`: {negative_prompt_2} and `negative_prompt_embeds`:" | |
f" {negative_prompt_embeds}. Please make sure to only forward one of the two." | |
) | |
if prompt_embeds is not None and negative_prompt_embeds is not None: | |
if prompt_embeds.shape != negative_prompt_embeds.shape: | |
raise ValueError( | |
"`prompt_embeds` and `negative_prompt_embeds` must have the same shape when passed directly, but" | |
f" got: `prompt_embeds` {prompt_embeds.shape} != `negative_prompt_embeds`" | |
f" {negative_prompt_embeds.shape}." | |
) | |
if padding_mask_crop is not None: | |
if not isinstance(image, PIL.Image.Image): | |
raise ValueError( | |
f"The image should be a PIL image when inpainting mask crop, but is of type" f" {type(image)}." | |
) | |
if not isinstance(mask_image, PIL.Image.Image): | |
raise ValueError( | |
f"The mask image should be a PIL image when inpainting mask crop, but is of type" | |
f" {type(mask_image)}." | |
) | |
if output_type != "pil": | |
raise ValueError(f"The output type should be PIL when inpainting mask crop, but is" f" {output_type}.") | |
def prepare_latents( | |
self, | |
batch_size, | |
num_channels_latents, | |
height, | |
width, | |
dtype, | |
device, | |
generator, | |
latents=None, | |
image=None, | |
timestep=None, | |
is_strength_max=True, | |
add_noise=True, | |
return_noise=False, | |
return_image_latents=False, | |
): | |
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor) | |
if isinstance(generator, list) and len(generator) != batch_size: | |
raise ValueError( | |
f"You have passed a list of generators of length {len(generator)}, but requested an effective batch" | |
f" size of {batch_size}. Make sure the batch size matches the length of the generators." | |
) | |
if (image is None or timestep is None) and not is_strength_max: | |
raise ValueError( | |
"Since strength < 1. initial latents are to be initialised as a combination of Image + Noise." | |
"However, either the image or the noise timestep has not been provided." | |
) | |
if image.shape[1] == 4: | |
image_latents = image.to(device=device, dtype=dtype) | |
image_latents = image_latents.repeat(batch_size // image_latents.shape[0], 1, 1, 1) | |
elif return_image_latents or (latents is None and not is_strength_max): | |
image = image.to(device=device, dtype=dtype) | |
image_latents = self._encode_vae_image(image=image, generator=generator) | |
image_latents = image_latents.repeat(batch_size // image_latents.shape[0], 1, 1, 1) | |
if latents is None and add_noise: | |
noise = randn_tensor(shape, generator=generator, device=device, dtype=dtype) | |
# if strength is 1. then initialise the latents to noise, else initial to image + noise | |
latents = noise if is_strength_max else self.scheduler.add_noise(image_latents, noise, timestep) | |
# if pure noise then scale the initial latents by the Scheduler's init sigma | |
latents = latents * self.scheduler.init_noise_sigma if is_strength_max else latents | |
elif add_noise: | |
noise = latents.to(device) | |
latents = noise * self.scheduler.init_noise_sigma | |
else: | |
noise = randn_tensor(shape, generator=generator, device=device, dtype=dtype) | |
latents = image_latents.to(device) | |
outputs = (latents,) | |
if return_noise: | |
outputs += (noise,) | |
if return_image_latents: | |
outputs += (image_latents,) | |
return outputs | |
def _encode_vae_image(self, image: torch.Tensor, generator: torch.Generator): | |
dtype = image.dtype | |
if self.vae.config.force_upcast: | |
image = image.float() | |
self.vae.to(dtype=torch.float32) | |
if isinstance(generator, list): | |
image_latents = [ | |
retrieve_latents(self.vae.encode(image[i : i + 1]), generator=generator[i]) | |
for i in range(image.shape[0]) | |
] | |
image_latents = torch.cat(image_latents, dim=0) | |
else: | |
image_latents = retrieve_latents(self.vae.encode(image), generator=generator) | |
if self.vae.config.force_upcast: | |
self.vae.to(dtype) | |
image_latents = image_latents.to(dtype) | |
image_latents = self.vae.config.scaling_factor * image_latents | |
return image_latents | |
def prepare_mask_latents( | |
self, mask, masked_image, batch_size, height, width, dtype, device, generator, do_classifier_free_guidance | |
): | |
# resize the mask to latents shape as we concatenate the mask to the latents | |
# we do that before converting to dtype to avoid breaking in case we're using cpu_offload | |
# and half precision | |
mask = torch.nn.functional.interpolate( | |
mask, size=(height // self.vae_scale_factor, width // self.vae_scale_factor) | |
) | |
mask = mask.to(device=device, dtype=dtype) | |
# duplicate mask and masked_image_latents for each generation per prompt, using mps friendly method | |
if mask.shape[0] < batch_size: | |
if not batch_size % mask.shape[0] == 0: | |
raise ValueError( | |
"The passed mask and the required batch size don't match. Masks are supposed to be duplicated to" | |
f" a total batch size of {batch_size}, but {mask.shape[0]} masks were passed. Make sure the number" | |
" of masks that you pass is divisible by the total requested batch size." | |
) | |
mask = mask.repeat(batch_size // mask.shape[0], 1, 1, 1) | |
mask = torch.cat([mask] * 2) if do_classifier_free_guidance else mask | |
if masked_image is not None and masked_image.shape[1] == 4: | |
masked_image_latents = masked_image | |
else: | |
masked_image_latents = None | |
if masked_image is not None: | |
if masked_image_latents is None: | |
masked_image = masked_image.to(device=device, dtype=dtype) | |
masked_image_latents = self._encode_vae_image(masked_image, generator=generator) | |
if masked_image_latents.shape[0] < batch_size: | |
if not batch_size % masked_image_latents.shape[0] == 0: | |
raise ValueError( | |
"The passed images and the required batch size don't match. Images are supposed to be duplicated" | |
f" to a total batch size of {batch_size}, but {masked_image_latents.shape[0]} images were passed." | |
" Make sure the number of images that you pass is divisible by the total requested batch size." | |
) | |
masked_image_latents = masked_image_latents.repeat( | |
batch_size // masked_image_latents.shape[0], 1, 1, 1 | |
) | |
masked_image_latents = ( | |
torch.cat([masked_image_latents] * 2) if do_classifier_free_guidance else masked_image_latents | |
) | |
# aligning device to prevent device errors when concating it with the latent model input | |
masked_image_latents = masked_image_latents.to(device=device, dtype=dtype) | |
return mask, masked_image_latents | |
# Copied from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl_img2img.StableDiffusionXLImg2ImgPipeline.get_timesteps | |
def get_timesteps(self, num_inference_steps, strength, device, denoising_start=None): | |
# get the original timestep using init_timestep | |
if denoising_start is None: | |
init_timestep = min(int(num_inference_steps * strength), num_inference_steps) | |
t_start = max(num_inference_steps - init_timestep, 0) | |
else: | |
t_start = 0 | |
timesteps = self.scheduler.timesteps[t_start * self.scheduler.order :] | |
# Strength is irrelevant if we directly request a timestep to start at; | |
# that is, strength is determined by the denoising_start instead. | |
if denoising_start is not None: | |
discrete_timestep_cutoff = int( | |
round( | |
self.scheduler.config.num_train_timesteps | |
- (denoising_start * self.scheduler.config.num_train_timesteps) | |
) | |
) | |
num_inference_steps = (timesteps < discrete_timestep_cutoff).sum().item() | |
if self.scheduler.order == 2 and num_inference_steps % 2 == 0: | |
# if the scheduler is a 2nd order scheduler we might have to do +1 | |
# because `num_inference_steps` might be even given that every timestep | |
# (except the highest one) is duplicated. If `num_inference_steps` is even it would | |
# mean that we cut the timesteps in the middle of the denoising step | |
# (between 1st and 2nd devirative) which leads to incorrect results. By adding 1 | |
# we ensure that the denoising process always ends after the 2nd derivate step of the scheduler | |
num_inference_steps = num_inference_steps + 1 | |
# because t_n+1 >= t_n, we slice the timesteps starting from the end | |
timesteps = timesteps[-num_inference_steps:] | |
return timesteps, num_inference_steps | |
return timesteps, num_inference_steps - t_start | |
# Copied from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl_img2img.StableDiffusionXLImg2ImgPipeline._get_add_time_ids | |
def _get_add_time_ids( | |
self, | |
original_size, | |
crops_coords_top_left, | |
target_size, | |
aesthetic_score, | |
negative_aesthetic_score, | |
negative_original_size, | |
negative_crops_coords_top_left, | |
negative_target_size, | |
dtype, | |
text_encoder_projection_dim=None, | |
): | |
if self.config.requires_aesthetics_score: | |
add_time_ids = list(original_size + crops_coords_top_left + (aesthetic_score,)) | |
add_neg_time_ids = list( | |
negative_original_size + negative_crops_coords_top_left + (negative_aesthetic_score,) | |
) | |
else: | |
add_time_ids = list(original_size + crops_coords_top_left + target_size) | |
add_neg_time_ids = list(negative_original_size + crops_coords_top_left + negative_target_size) | |
passed_add_embed_dim = ( | |
self.unet.config.addition_time_embed_dim * len(add_time_ids) + text_encoder_projection_dim | |
) | |
expected_add_embed_dim = self.unet.add_embedding.linear_1.in_features | |
if ( | |
expected_add_embed_dim > passed_add_embed_dim | |
and (expected_add_embed_dim - passed_add_embed_dim) == self.unet.config.addition_time_embed_dim | |
): | |
raise ValueError( | |
f"Model expects an added time embedding vector of length {expected_add_embed_dim}, but a vector of {passed_add_embed_dim} was created. Please make sure to enable `requires_aesthetics_score` with `pipe.register_to_config(requires_aesthetics_score=True)` to make sure `aesthetic_score` {aesthetic_score} and `negative_aesthetic_score` {negative_aesthetic_score} is correctly used by the model." | |
) | |
elif ( | |
expected_add_embed_dim < passed_add_embed_dim | |
and (passed_add_embed_dim - expected_add_embed_dim) == self.unet.config.addition_time_embed_dim | |
): | |
raise ValueError( | |
f"Model expects an added time embedding vector of length {expected_add_embed_dim}, but a vector of {passed_add_embed_dim} was created. Please make sure to disable `requires_aesthetics_score` with `pipe.register_to_config(requires_aesthetics_score=False)` to make sure `target_size` {target_size} is correctly used by the model." | |
) | |
elif expected_add_embed_dim != passed_add_embed_dim: | |
raise ValueError( | |
f"Model expects an added time embedding vector of length {expected_add_embed_dim}, but a vector of {passed_add_embed_dim} was created. The model has an incorrect config. Please check `unet.config.time_embedding_type` and `text_encoder_2.config.projection_dim`." | |
) | |
add_time_ids = torch.tensor([add_time_ids], dtype=dtype) | |
add_neg_time_ids = torch.tensor([add_neg_time_ids], dtype=dtype) | |
return add_time_ids, add_neg_time_ids | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_upscale.StableDiffusionUpscalePipeline.upcast_vae | |
def upcast_vae(self): | |
dtype = self.vae.dtype | |
self.vae.to(dtype=torch.float32) | |
use_torch_2_0_or_xformers = isinstance( | |
self.vae.decoder.mid_block.attentions[0].processor, | |
( | |
AttnProcessor2_0, | |
XFormersAttnProcessor, | |
LoRAXFormersAttnProcessor, | |
LoRAAttnProcessor2_0, | |
), | |
) | |
# if xformers or torch_2_0 is used attention block does not need | |
# to be in float32 which can save lots of memory | |
if use_torch_2_0_or_xformers: | |
self.vae.post_quant_conv.to(dtype) | |
self.vae.decoder.conv_in.to(dtype) | |
self.vae.decoder.mid_block.to(dtype) | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_freeu | |
def enable_freeu(self, s1: float, s2: float, b1: float, b2: float): | |
r"""Enables the FreeU mechanism as in https://arxiv.org/abs/2309.11497. | |
The suffixes after the scaling factors represent the stages where they are being applied. | |
Please refer to the [official repository](https://github.com/ChenyangSi/FreeU) for combinations of the values | |
that are known to work well for different pipelines such as Stable Diffusion v1, v2, and Stable Diffusion XL. | |
Args: | |
s1 (`float`): | |
Scaling factor for stage 1 to attenuate the contributions of the skip features. This is done to | |
mitigate "oversmoothing effect" in the enhanced denoising process. | |
s2 (`float`): | |
Scaling factor for stage 2 to attenuate the contributions of the skip features. This is done to | |
mitigate "oversmoothing effect" in the enhanced denoising process. | |
b1 (`float`): Scaling factor for stage 1 to amplify the contributions of backbone features. | |
b2 (`float`): Scaling factor for stage 2 to amplify the contributions of backbone features. | |
""" | |
if not hasattr(self, "unet"): | |
raise ValueError("The pipeline must have `unet` for using FreeU.") | |
self.unet.enable_freeu(s1=s1, s2=s2, b1=b1, b2=b2) | |
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_freeu | |
def disable_freeu(self): | |
"""Disables the FreeU mechanism if enabled.""" | |
self.unet.disable_freeu() | |
# Copied from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl.StableDiffusionXLPipeline.fuse_qkv_projections | |
def fuse_qkv_projections(self, unet: bool = True, vae: bool = True): | |
""" | |
Enables fused QKV projections. For self-attention modules, all projection matrices (i.e., query, | |
key, value) are fused. For cross-attention modules, key and value projection matrices are fused. | |
<Tip warning={true}> | |
This API is π§ͺ experimental. | |
</Tip> | |
Args: | |
unet (`bool`, defaults to `True`): To apply fusion on the UNet. | |
vae (`bool`, defaults to `True`): To apply fusion on the VAE. | |
""" | |
self.fusing_unet = False | |
self.fusing_vae = False | |
if unet: | |
self.fusing_unet = True | |
self.unet.fuse_qkv_projections() | |
self.unet.set_attn_processor(FusedAttnProcessor2_0()) | |
if vae: | |
if not isinstance(self.vae, AutoencoderKL): | |
raise ValueError("`fuse_qkv_projections()` is only supported for the VAE of type `AutoencoderKL`.") | |
self.fusing_vae = True | |
self.vae.fuse_qkv_projections() | |
self.vae.set_attn_processor(FusedAttnProcessor2_0()) | |
# Copied from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl.StableDiffusionXLPipeline.unfuse_qkv_projections | |
def unfuse_qkv_projections(self, unet: bool = True, vae: bool = True): | |
"""Disable QKV projection fusion if enabled. | |
<Tip warning={true}> | |
This API is π§ͺ experimental. | |
</Tip> | |
Args: | |
unet (`bool`, defaults to `True`): To apply fusion on the UNet. | |
vae (`bool`, defaults to `True`): To apply fusion on the VAE. | |
""" | |
if unet: | |
if not self.fusing_unet: | |
logger.warning("The UNet was not initially fused for QKV projections. Doing nothing.") | |
else: | |
self.unet.unfuse_qkv_projections() | |
self.fusing_unet = False | |
if vae: | |
if not self.fusing_vae: | |
logger.warning("The VAE was not initially fused for QKV projections. Doing nothing.") | |
else: | |
self.vae.unfuse_qkv_projections() | |
self.fusing_vae = False | |
# Copied from diffusers.pipelines.latent_consistency_models.pipeline_latent_consistency_text2img.LatentConsistencyModelPipeline.get_guidance_scale_embedding | |
def get_guidance_scale_embedding(self, w, embedding_dim=512, dtype=torch.float32): | |
""" | |
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298 | |
Args: | |
timesteps (`torch.Tensor`): | |
generate embedding vectors at these timesteps | |
embedding_dim (`int`, *optional*, defaults to 512): | |
dimension of the embeddings to generate | |
dtype: | |
data type of the generated embeddings | |
Returns: | |
`torch.FloatTensor`: Embedding vectors with shape `(len(timesteps), embedding_dim)` | |
""" | |
assert len(w.shape) == 1 | |
w = w * 1000.0 | |
half_dim = embedding_dim // 2 | |
emb = torch.log(torch.tensor(10000.0)) / (half_dim - 1) | |
emb = torch.exp(torch.arange(half_dim, dtype=dtype) * -emb) | |
emb = w.to(dtype)[:, None] * emb[None, :] | |
emb = torch.cat([torch.sin(emb), torch.cos(emb)], dim=1) | |
if embedding_dim % 2 == 1: # zero pad | |
emb = torch.nn.functional.pad(emb, (0, 1)) | |
assert emb.shape == (w.shape[0], embedding_dim) | |
return emb | |
def guidance_scale(self): | |
return self._guidance_scale | |
def guidance_rescale(self): | |
return self._guidance_rescale | |
def clip_skip(self): | |
return self._clip_skip | |
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2) | |
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1` | |
# corresponds to doing no classifier free guidance. | |
def do_classifier_free_guidance(self): | |
return self._guidance_scale > 1 and self.unet.config.time_cond_proj_dim is None | |
def cross_attention_kwargs(self): | |
return self._cross_attention_kwargs | |
def denoising_end(self): | |
return self._denoising_end | |
def denoising_start(self): | |
return self._denoising_start | |
def num_timesteps(self): | |
return self._num_timesteps | |
def interrupt(self): | |
return self._interrupt | |
def __call__( | |
self, | |
prompt: Union[str, List[str]] = None, | |
prompt_2: Optional[Union[str, List[str]]] = None, | |
image: PipelineImageInput = None, | |
mask_image: PipelineImageInput = None, | |
masked_image_latents: torch.FloatTensor = None, | |
height: Optional[int] = None, | |
width: Optional[int] = None, | |
padding_mask_crop: Optional[int] = None, | |
strength: float = 0.9999, | |
num_inference_steps: int = 50, | |
timesteps: List[int] = None, | |
denoising_start: Optional[float] = None, | |
denoising_end: Optional[float] = None, | |
guidance_scale: float = 7.5, | |
negative_prompt: Optional[Union[str, List[str]]] = None, | |
negative_prompt_2: Optional[Union[str, List[str]]] = None, | |
num_images_per_prompt: Optional[int] = 1, | |
eta: float = 0.0, | |
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None, | |
latents: Optional[torch.FloatTensor] = None, | |
prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
negative_pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
ip_adapter_image: Optional[PipelineImageInput] = None, | |
output_type: Optional[str] = "pil", | |
cloth =None, | |
pose_img = None, | |
text_embeds_cloth=None, | |
return_dict: bool = True, | |
cross_attention_kwargs: Optional[Dict[str, Any]] = None, | |
guidance_rescale: float = 0.0, | |
original_size: Tuple[int, int] = None, | |
crops_coords_top_left: Tuple[int, int] = (0, 0), | |
target_size: Tuple[int, int] = None, | |
negative_original_size: Optional[Tuple[int, int]] = None, | |
negative_crops_coords_top_left: Tuple[int, int] = (0, 0), | |
negative_target_size: Optional[Tuple[int, int]] = None, | |
aesthetic_score: float = 6.0, | |
negative_aesthetic_score: float = 2.5, | |
clip_skip: Optional[int] = None, | |
pooled_prompt_embeds_c=None, | |
callback_on_step_end: Optional[Callable[[int, int, Dict], None]] = None, | |
callback_on_step_end_tensor_inputs: List[str] = ["latents"], | |
**kwargs, | |
): | |
r""" | |
Function invoked when calling the pipeline for generation. | |
Args: | |
prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`. | |
instead. | |
prompt_2 (`str` or `List[str]`, *optional*): | |
The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is | |
used in both text-encoders | |
image (`PIL.Image.Image`): | |
`Image`, or tensor representing an image batch which will be inpainted, *i.e.* parts of the image will | |
be masked out with `mask_image` and repainted according to `prompt`. | |
mask_image (`PIL.Image.Image`): | |
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be | |
repainted, while black pixels will be preserved. If `mask_image` is a PIL image, it will be converted | |
to a single channel (luminance) before use. If it's a tensor, it should contain one color channel (L) | |
instead of 3, so the expected shape would be `(B, H, W, 1)`. | |
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor): | |
The height in pixels of the generated image. This is set to 1024 by default for the best results. | |
Anything below 512 pixels won't work well for | |
[stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) | |
and checkpoints that are not specifically fine-tuned on low resolutions. | |
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor): | |
The width in pixels of the generated image. This is set to 1024 by default for the best results. | |
Anything below 512 pixels won't work well for | |
[stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) | |
and checkpoints that are not specifically fine-tuned on low resolutions. | |
padding_mask_crop (`int`, *optional*, defaults to `None`): | |
The size of margin in the crop to be applied to the image and masking. If `None`, no crop is applied to image and mask_image. If | |
`padding_mask_crop` is not `None`, it will first find a rectangular region with the same aspect ration of the image and | |
contains all masked area, and then expand that area based on `padding_mask_crop`. The image and mask_image will then be cropped based on | |
the expanded area before resizing to the original image size for inpainting. This is useful when the masked area is small while the image is large | |
and contain information inreleant for inpainging, such as background. | |
strength (`float`, *optional*, defaults to 0.9999): | |
Conceptually, indicates how much to transform the masked portion of the reference `image`. Must be | |
between 0 and 1. `image` will be used as a starting point, adding more noise to it the larger the | |
`strength`. The number of denoising steps depends on the amount of noise initially added. When | |
`strength` is 1, added noise will be maximum and the denoising process will run for the full number of | |
iterations specified in `num_inference_steps`. A value of 1, therefore, essentially ignores the masked | |
portion of the reference `image`. Note that in the case of `denoising_start` being declared as an | |
integer, the value of `strength` will be ignored. | |
num_inference_steps (`int`, *optional*, defaults to 50): | |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the | |
expense of slower inference. | |
timesteps (`List[int]`, *optional*): | |
Custom timesteps to use for the denoising process with schedulers which support a `timesteps` argument | |
in their `set_timesteps` method. If not defined, the default behavior when `num_inference_steps` is | |
passed will be used. Must be in descending order. | |
denoising_start (`float`, *optional*): | |
When specified, indicates the fraction (between 0.0 and 1.0) of the total denoising process to be | |
bypassed before it is initiated. Consequently, the initial part of the denoising process is skipped and | |
it is assumed that the passed `image` is a partly denoised image. Note that when this is specified, | |
strength will be ignored. The `denoising_start` parameter is particularly beneficial when this pipeline | |
is integrated into a "Mixture of Denoisers" multi-pipeline setup, as detailed in [**Refining the Image | |
Output**](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/stable_diffusion_xl#refining-the-image-output). | |
denoising_end (`float`, *optional*): | |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be | |
completed before it is intentionally prematurely terminated. As a result, the returned sample will | |
still retain a substantial amount of noise (ca. final 20% of timesteps still needed) and should be | |
denoised by a successor pipeline that has `denoising_start` set to 0.8 so that it only denoises the | |
final 20% of the scheduler. The denoising_end parameter should ideally be utilized when this pipeline | |
forms a part of a "Mixture of Denoisers" multi-pipeline setup, as elaborated in [**Refining the Image | |
Output**](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/stable_diffusion_xl#refining-the-image-output). | |
guidance_scale (`float`, *optional*, defaults to 7.5): | |
Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598). | |
`guidance_scale` is defined as `w` of equation 2. of [Imagen | |
Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale > | |
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`, | |
usually at the expense of lower image quality. | |
negative_prompt (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
`negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
less than `1`). | |
negative_prompt_2 (`str` or `List[str]`, *optional*): | |
The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and | |
`text_encoder_2`. If not defined, `negative_prompt` is used in both text-encoders | |
prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
provided, text embeddings will be generated from `prompt` input argument. | |
negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
argument. | |
pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. | |
If not provided, pooled text embeddings will be generated from `prompt` input argument. | |
negative_pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` | |
input argument. | |
ip_adapter_image: (`PipelineImageInput`, *optional*): Optional image input to work with IP Adapters. | |
num_images_per_prompt (`int`, *optional*, defaults to 1): | |
The number of images to generate per prompt. | |
eta (`float`, *optional*, defaults to 0.0): | |
Corresponds to parameter eta (Ξ·) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to | |
[`schedulers.DDIMScheduler`], will be ignored for others. | |
generator (`torch.Generator`, *optional*): | |
One or a list of [torch generator(s)](https://pytorch.org/docs/stable/generated/torch.Generator.html) | |
to make generation deterministic. | |
latents (`torch.FloatTensor`, *optional*): | |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image | |
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents | |
tensor will ge generated by sampling using the supplied random `generator`. | |
output_type (`str`, *optional*, defaults to `"pil"`): | |
The output format of the generate image. Choose between | |
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`. | |
return_dict (`bool`, *optional*, defaults to `True`): | |
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a | |
plain tuple. | |
cross_attention_kwargs (`dict`, *optional*): | |
A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under | |
`self.processor` in | |
[diffusers.models.attention_processor](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py). | |
original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
If `original_size` is not the same as `target_size` the image will appear to be down- or upsampled. | |
`original_size` defaults to `(height, width)` if not specified. Part of SDXL's micro-conditioning as | |
explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): | |
`crops_coords_top_left` can be used to generate an image that appears to be "cropped" from the position | |
`crops_coords_top_left` downwards. Favorable, well-centered images are usually achieved by setting | |
`crops_coords_top_left` to (0, 0). Part of SDXL's micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
For most cases, `target_size` should be set to the desired height and width of the generated image. If | |
not specified it will default to `(height, width)`. Part of SDXL's micro-conditioning as explained in | |
section 2.2 of [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
negative_original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's | |
micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more | |
information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. | |
negative_crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): | |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's | |
micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more | |
information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. | |
negative_target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
To negatively condition the generation process based on a target image resolution. It should be as same | |
as the `target_size` for most cases. Part of SDXL's micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). For more | |
information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. | |
aesthetic_score (`float`, *optional*, defaults to 6.0): | |
Used to simulate an aesthetic score of the generated image by influencing the positive text condition. | |
Part of SDXL's micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
negative_aesthetic_score (`float`, *optional*, defaults to 2.5): | |
Part of SDXL's micro-conditioning as explained in section 2.2 of | |
[https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). Can be used to | |
simulate an aesthetic score of the generated image by influencing the negative text condition. | |
clip_skip (`int`, *optional*): | |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that | |
the output of the pre-final layer will be used for computing the prompt embeddings. | |
callback_on_step_end (`Callable`, *optional*): | |
A function that calls at the end of each denoising steps during the inference. The function is called | |
with the following arguments: `callback_on_step_end(self: DiffusionPipeline, step: int, timestep: int, | |
callback_kwargs: Dict)`. `callback_kwargs` will include a list of all tensors as specified by | |
`callback_on_step_end_tensor_inputs`. | |
callback_on_step_end_tensor_inputs (`List`, *optional*): | |
The list of tensor inputs for the `callback_on_step_end` function. The tensors specified in the list | |
will be passed as `callback_kwargs` argument. You will only be able to include variables listed in the | |
`._callback_tensor_inputs` attribute of your pipeline class. | |
Examples: | |
Returns: | |
[`~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput`] or `tuple`: | |
[`~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput`] if `return_dict` is True, otherwise a | |
`tuple. `tuple. When returning a tuple, the first element is a list with the generated images. | |
""" | |
callback = kwargs.pop("callback", None) | |
callback_steps = kwargs.pop("callback_steps", None) | |
if callback is not None: | |
deprecate( | |
"callback", | |
"1.0.0", | |
"Passing `callback` as an input argument to `__call__` is deprecated, consider use `callback_on_step_end`", | |
) | |
if callback_steps is not None: | |
deprecate( | |
"callback_steps", | |
"1.0.0", | |
"Passing `callback_steps` as an input argument to `__call__` is deprecated, consider use `callback_on_step_end`", | |
) | |
# 0. Default height and width to unet | |
height = height or self.unet.config.sample_size * self.vae_scale_factor | |
width = width or self.unet.config.sample_size * self.vae_scale_factor | |
# 1. Check inputs | |
self.check_inputs( | |
prompt, | |
prompt_2, | |
image, | |
mask_image, | |
height, | |
width, | |
strength, | |
callback_steps, | |
output_type, | |
negative_prompt, | |
negative_prompt_2, | |
prompt_embeds, | |
negative_prompt_embeds, | |
callback_on_step_end_tensor_inputs, | |
padding_mask_crop, | |
) | |
self._guidance_scale = guidance_scale | |
self._guidance_rescale = guidance_rescale | |
self._clip_skip = clip_skip | |
self._cross_attention_kwargs = cross_attention_kwargs | |
self._denoising_end = denoising_end | |
self._denoising_start = denoising_start | |
self._interrupt = False | |
# 2. Define call parameters | |
if prompt is not None and isinstance(prompt, str): | |
batch_size = 1 | |
elif prompt is not None and isinstance(prompt, list): | |
batch_size = len(prompt) | |
else: | |
batch_size = prompt_embeds.shape[0] | |
device = self._execution_device | |
# 3. Encode input prompt | |
text_encoder_lora_scale = ( | |
self.cross_attention_kwargs.get("scale", None) if self.cross_attention_kwargs is not None else None | |
) | |
( | |
prompt_embeds, | |
negative_prompt_embeds, | |
pooled_prompt_embeds, | |
negative_pooled_prompt_embeds, | |
) = self.encode_prompt( | |
prompt=prompt, | |
prompt_2=prompt_2, | |
device=device, | |
num_images_per_prompt=num_images_per_prompt, | |
do_classifier_free_guidance=self.do_classifier_free_guidance, | |
negative_prompt=negative_prompt, | |
negative_prompt_2=negative_prompt_2, | |
prompt_embeds=prompt_embeds, | |
negative_prompt_embeds=negative_prompt_embeds, | |
pooled_prompt_embeds=pooled_prompt_embeds, | |
negative_pooled_prompt_embeds=negative_pooled_prompt_embeds, | |
lora_scale=text_encoder_lora_scale, | |
clip_skip=self.clip_skip, | |
) | |
# 4. set timesteps | |
def denoising_value_valid(dnv): | |
return isinstance(self.denoising_end, float) and 0 < dnv < 1 | |
timesteps, num_inference_steps = retrieve_timesteps(self.scheduler, num_inference_steps, device, timesteps) | |
timesteps, num_inference_steps = self.get_timesteps( | |
num_inference_steps, | |
strength, | |
device, | |
denoising_start=self.denoising_start if denoising_value_valid else None, | |
) | |
# check that number of inference steps is not < 1 - as this doesn't make sense | |
if num_inference_steps < 1: | |
raise ValueError( | |
f"After adjusting the num_inference_steps by strength parameter: {strength}, the number of pipeline" | |
f"steps is {num_inference_steps} which is < 1 and not appropriate for this pipeline." | |
) | |
# at which timestep to set the initial noise (n.b. 50% if strength is 0.5) | |
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt) | |
# create a boolean to check if the strength is set to 1. if so then initialise the latents with pure noise | |
is_strength_max = strength == 1.0 | |
# 5. Preprocess mask and image | |
if padding_mask_crop is not None: | |
crops_coords = self.mask_processor.get_crop_region(mask_image, width, height, pad=padding_mask_crop) | |
resize_mode = "fill" | |
else: | |
crops_coords = None | |
resize_mode = "default" | |
original_image = image | |
init_image = self.image_processor.preprocess( | |
image, height=height, width=width, crops_coords=crops_coords, resize_mode=resize_mode | |
) | |
init_image = init_image.to(dtype=torch.float32) | |
mask = self.mask_processor.preprocess( | |
mask_image, height=height, width=width, resize_mode=resize_mode, crops_coords=crops_coords | |
) | |
if masked_image_latents is not None: | |
masked_image = masked_image_latents | |
elif init_image.shape[1] == 4: | |
# if images are in latent space, we can't mask it | |
masked_image = None | |
else: | |
masked_image = init_image * (mask < 0.5) | |
# 6. Prepare latent variables | |
num_channels_latents = self.vae.config.latent_channels | |
num_channels_unet = self.unet.config.in_channels | |
return_image_latents = num_channels_unet == 4 | |
add_noise = True if self.denoising_start is None else False | |
latents_outputs = self.prepare_latents( | |
batch_size * num_images_per_prompt, | |
num_channels_latents, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
latents, | |
image=init_image, | |
timestep=latent_timestep, | |
is_strength_max=is_strength_max, | |
add_noise=add_noise, | |
return_noise=True, | |
return_image_latents=return_image_latents, | |
) | |
if return_image_latents: | |
latents, noise, image_latents = latents_outputs | |
else: | |
latents, noise = latents_outputs | |
# 7. Prepare mask latent variables | |
mask, masked_image_latents = self.prepare_mask_latents( | |
mask, | |
masked_image, | |
batch_size * num_images_per_prompt, | |
height, | |
width, | |
prompt_embeds.dtype, | |
device, | |
generator, | |
self.do_classifier_free_guidance, | |
) | |
pose_img = pose_img.to(device=device, dtype=prompt_embeds.dtype) | |
pose_img = self.vae.encode(pose_img).latent_dist.sample() | |
pose_img = pose_img * self.vae.config.scaling_factor | |
# pose_img = self._encode_vae_image(pose_img, generator=generator) | |
pose_img = ( | |
torch.cat([pose_img] * 2) if self.do_classifier_free_guidance else pose_img | |
) | |
cloth = self._encode_vae_image(cloth, generator=generator) | |
# # 8. Check that sizes of mask, masked image and latents match | |
# if num_channels_unet == 9: | |
# # default case for runwayml/stable-diffusion-inpainting | |
# num_channels_mask = mask.shape[1] | |
# num_channels_masked_image = masked_image_latents.shape[1] | |
# if num_channels_latents + num_channels_mask + num_channels_masked_image != self.unet.config.in_channels: | |
# raise ValueError( | |
# f"Incorrect configuration settings! The config of `pipeline.unet`: {self.unet.config} expects" | |
# f" {self.unet.config.in_channels} but received `num_channels_latents`: {num_channels_latents} +" | |
# f" `num_channels_mask`: {num_channels_mask} + `num_channels_masked_image`: {num_channels_masked_image}" | |
# f" = {num_channels_latents+num_channels_masked_image+num_channels_mask}. Please verify the config of" | |
# " `pipeline.unet` or your `mask_image` or `image` input." | |
# ) | |
# elif num_channels_unet != 4: | |
# raise ValueError( | |
# f"The unet {self.unet.__class__} should have either 4 or 9 input channels, not {self.unet.config.in_channels}." | |
# ) | |
# 8.1 Prepare extra step kwargs. | |
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta) | |
# 9. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline | |
height, width = latents.shape[-2:] | |
height = height * self.vae_scale_factor | |
width = width * self.vae_scale_factor | |
original_size = original_size or (height, width) | |
target_size = target_size or (height, width) | |
# 10. Prepare added time ids & embeddings | |
if negative_original_size is None: | |
negative_original_size = original_size | |
if negative_target_size is None: | |
negative_target_size = target_size | |
add_text_embeds = pooled_prompt_embeds | |
if self.text_encoder_2 is None: | |
text_encoder_projection_dim = int(pooled_prompt_embeds.shape[-1]) | |
else: | |
text_encoder_projection_dim = self.text_encoder_2.config.projection_dim | |
add_time_ids, add_neg_time_ids = self._get_add_time_ids( | |
original_size, | |
crops_coords_top_left, | |
target_size, | |
aesthetic_score, | |
negative_aesthetic_score, | |
negative_original_size, | |
negative_crops_coords_top_left, | |
negative_target_size, | |
dtype=prompt_embeds.dtype, | |
text_encoder_projection_dim=text_encoder_projection_dim, | |
) | |
add_time_ids = add_time_ids.repeat(batch_size * num_images_per_prompt, 1) | |
if self.do_classifier_free_guidance: | |
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds], dim=0) | |
add_text_embeds = torch.cat([negative_pooled_prompt_embeds, add_text_embeds], dim=0) | |
add_neg_time_ids = add_neg_time_ids.repeat(batch_size * num_images_per_prompt, 1) | |
add_time_ids = torch.cat([add_neg_time_ids, add_time_ids], dim=0) | |
prompt_embeds = prompt_embeds.to(device) | |
add_text_embeds = add_text_embeds.to(device) | |
add_time_ids = add_time_ids.to(device) | |
if ip_adapter_image is not None: | |
image_embeds = self.prepare_ip_adapter_image_embeds( | |
ip_adapter_image, device, batch_size * num_images_per_prompt | |
) | |
#project outside for loop | |
image_embeds = unet.encoder_hid_proj(image_embeds).to(prompt_embeds.dtype) | |
# 11. Denoising loop | |
num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0) | |
if ( | |
self.denoising_end is not None | |
and self.denoising_start is not None | |
and denoising_value_valid(self.denoising_end) | |
and denoising_value_valid(self.denoising_start) | |
and self.denoising_start >= self.denoising_end | |
): | |
raise ValueError( | |
f"`denoising_start`: {self.denoising_start} cannot be larger than or equal to `denoising_end`: " | |
+ f" {self.denoising_end} when using type float." | |
) | |
elif self.denoising_end is not None and denoising_value_valid(self.denoising_end): | |
discrete_timestep_cutoff = int( | |
round( | |
self.scheduler.config.num_train_timesteps | |
- (self.denoising_end * self.scheduler.config.num_train_timesteps) | |
) | |
) | |
num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps))) | |
timesteps = timesteps[:num_inference_steps] | |
# 11.1 Optionally get Guidance Scale Embedding | |
timestep_cond = None | |
if self.unet.config.time_cond_proj_dim is not None: | |
guidance_scale_tensor = torch.tensor(self.guidance_scale - 1).repeat(batch_size * num_images_per_prompt) | |
timestep_cond = self.get_guidance_scale_embedding( | |
guidance_scale_tensor, embedding_dim=self.unet.config.time_cond_proj_dim | |
).to(device=device, dtype=latents.dtype) | |
self._num_timesteps = len(timesteps) | |
with self.progress_bar(total=num_inference_steps) as progress_bar: | |
for i, t in enumerate(timesteps): | |
if self.interrupt: | |
continue | |
# expand the latents if we are doing classifier free guidance | |
latent_model_input = torch.cat([latents] * 2) if self.do_classifier_free_guidance else latents | |
# concat latents, mask, masked_image_latents in the channel dimension | |
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t) | |
# bsz = mask.shape[0] | |
if num_channels_unet == 13: | |
latent_model_input = torch.cat([latent_model_input, mask, masked_image_latents,pose_img], dim=1) | |
# if num_channels_unet == 9: | |
# latent_model_input = torch.cat([latent_model_input, mask, masked_image_latents], dim=1) | |
# predict the noise residual | |
added_cond_kwargs = {"text_embeds": add_text_embeds, "time_ids": add_time_ids} | |
if ip_adapter_image is not None: | |
added_cond_kwargs["image_embeds"] = image_embeds | |
# down,reference_features = self.UNet_Encoder(cloth,t, text_embeds_cloth,added_cond_kwargs= {"text_embeds": pooled_prompt_embeds_c, "time_ids": add_time_ids},return_dict=False) | |
down,reference_features = self.unet_encoder(cloth,t, text_embeds_cloth,return_dict=False) | |
# print(type(reference_features)) | |
# print(reference_features) | |
reference_features = list(reference_features) | |
# print(len(reference_features)) | |
# for elem in reference_features: | |
# print(elem.shape) | |
# exit(1) | |
if self.do_classifier_free_guidance: | |
reference_features = [torch.cat([torch.zeros_like(d), d]) for d in reference_features] | |
noise_pred = self.unet( | |
latent_model_input, | |
t, | |
encoder_hidden_states=prompt_embeds, | |
timestep_cond=timestep_cond, | |
cross_attention_kwargs=self.cross_attention_kwargs, | |
added_cond_kwargs=added_cond_kwargs, | |
return_dict=False, | |
garment_features=reference_features, | |
)[0] | |
# noise_pred = self.unet(latent_model_input, t, | |
# prompt_embeds,timestep_cond=timestep_cond,cross_attention_kwargs=self.cross_attention_kwargs,added_cond_kwargs=added_cond_kwargs,down_block_additional_attn=down ).sample | |
# perform guidance | |
if self.do_classifier_free_guidance: | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + self.guidance_scale * (noise_pred_text - noise_pred_uncond) | |
if self.do_classifier_free_guidance and self.guidance_rescale > 0.0: | |
# Based on 3.4. in https://arxiv.org/pdf/2305.08891.pdf | |
noise_pred = rescale_noise_cfg(noise_pred, noise_pred_text, guidance_rescale=self.guidance_rescale) | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0] | |
if num_channels_unet == 4: | |
init_latents_proper = image_latents | |
if self.do_classifier_free_guidance: | |
init_mask, _ = mask.chunk(2) | |
else: | |
init_mask = mask | |
if i < len(timesteps) - 1: | |
noise_timestep = timesteps[i + 1] | |
init_latents_proper = self.scheduler.add_noise( | |
init_latents_proper, noise, torch.tensor([noise_timestep]) | |
) | |
latents = (1 - init_mask) * init_latents_proper + init_mask * latents | |
if callback_on_step_end is not None: | |
callback_kwargs = {} | |
for k in callback_on_step_end_tensor_inputs: | |
callback_kwargs[k] = locals()[k] | |
callback_outputs = callback_on_step_end(self, i, t, callback_kwargs) | |
latents = callback_outputs.pop("latents", latents) | |
prompt_embeds = callback_outputs.pop("prompt_embeds", prompt_embeds) | |
negative_prompt_embeds = callback_outputs.pop("negative_prompt_embeds", negative_prompt_embeds) | |
add_text_embeds = callback_outputs.pop("add_text_embeds", add_text_embeds) | |
negative_pooled_prompt_embeds = callback_outputs.pop( | |
"negative_pooled_prompt_embeds", negative_pooled_prompt_embeds | |
) | |
add_time_ids = callback_outputs.pop("add_time_ids", add_time_ids) | |
add_neg_time_ids = callback_outputs.pop("add_neg_time_ids", add_neg_time_ids) | |
mask = callback_outputs.pop("mask", mask) | |
masked_image_latents = callback_outputs.pop("masked_image_latents", masked_image_latents) | |
# call the callback, if provided | |
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0): | |
progress_bar.update() | |
if callback is not None and i % callback_steps == 0: | |
step_idx = i // getattr(self.scheduler, "order", 1) | |
callback(step_idx, t, latents) | |
if XLA_AVAILABLE: | |
xm.mark_step() | |
if not output_type == "latent": | |
# make sure the VAE is in float32 mode, as it overflows in float16 | |
needs_upcasting = self.vae.dtype == torch.float16 and self.vae.config.force_upcast | |
if needs_upcasting: | |
self.upcast_vae() | |
latents = latents.to(next(iter(self.vae.post_quant_conv.parameters())).dtype) | |
image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0] | |
# cast back to fp16 if needed | |
if needs_upcasting: | |
self.vae.to(dtype=torch.float16) | |
# else: | |
# return StableDiffusionXLPipelineOutput(images=latents) | |
image = self.image_processor.postprocess(image, output_type=output_type) | |
if padding_mask_crop is not None: | |
image = [self.image_processor.apply_overlay(mask_image, original_image, i, crops_coords) for i in image] | |
# Offload all models | |
self.maybe_free_model_hooks() | |
# if not return_dict: | |
return (image,) | |
# return StableDiffusionXLPipelineOutput(images=image) |